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Set Definitions
Elements are distinct and sequence-independent ∣X∣=cardinality of X=number of elements in X
| Symbol | Set | Examples |
| | | |
| R | Reals | π, 2, 21 |
| Q | Rationals | 2,21 |
| Z | Integers | -1, 0, 1 |
| N | Naturals | 1, 2, 3 | Z− represents the negative integers Z+ represents the positive integers ∅ denotes the null or empty set A⊂B means A is a subset of B A∪B is the union of A and B A∩B is the intersection of A and B A′ or Aˉ or AC is the complement of A
Commutative: A∪B=B∪A,A∩B=B∩A
Associative: (A∪B)∪C=A∪(B∪C),(A∩B)∩C=A∩(B∩C)
Distributive: A∩(B∪C)=(A∩B)∪(A∩C),A∪(B∩C)=(A∪B)∩(A∪C)
De Morgan’s Laws: (A∪B)′=A′∩B′,(A∩B)′=A′∪B′
Cartesian Product
E.g.: Z×Z can represent the xy-plane
Injective (one-to-one): ∀x1,x2∈X,f(x1)=f(x2)→x1=x2
Each element in Y has either one or zero elements that map to it from X f passes the horizontal line test
Surjective (onto): ∀y∈Y,∃x∈X:f(x)=y
Each element in Y has at least one element that maps to it from X Ran(f)=Y Bijective (one-to-one correspondence): Injective and Surjective Relation: aRb=(a,b)∈R Reflexive: ∀x∈X,xRx Symmetric: ∀x,y∈X,xRy→yRx Anti-symmetric: ∀x,y∈X,xRy∧yRx→x=y Transitive: ∀x,y,z∈X,xRy∧yRz→xRz Partial Order: reflexive, anti-symmetric, and transitive Equivalence Relation: reflexive, symmetric, and transitive ∣X∣=n→∣P(X)∣=2n X∪Y=∅→∣X∪Y∣=∣X∣+∣Y∣ (Addition Principle) ∣X∪Y∣=∣X∣+∣Y∣−∣X∩Y∣
Logic ∴
Counting 🔢
P(n,k)=nPk=k!(kn)=(n−k)!n!
The number of permutations (these are ordered combination basically) when picking k elements from a set of size n is P(n,k)
E.g.: There is a set X={A,B,C,D}. How many ways are there to pick two elements from X when ordering matters? P(4,2)=12
E.g.: There are 4 empty parking spots. How many ways are there to park 2 unique cars when ordering matters? P(4,2)=12
C(n,k)=nCk=(kn)=k!P(n,r)=k!(n−k)!n!
Same as permutations, but ordering doesn’t matter
For the first example for P(n,k), there would be (24)=6 ways to pick two elements from X when ordering doesn’t matter
E.g.: How many ways are there to park 2 identical cars in 4 parking spots? (24)=6 (kn)=(n−kn),(kn+k)=(nn+k)
| | Repetition/replacement allowed | Repetition/replacement not allowed |
| | | |
| Order matters | nk | P(n,k) |
| Order doesn’t matter | (kn+k−1) | C(n,k) |
(kn+k−1)=(n−1n+k−1)
This is the number of ways to uniquely choose k items from n types of items when ordering doesn’t matter.
This is “with repetition” because you are allowed to pick multiple of the same type
E.g.: What is the number of ways to choose nine pieces of bread when there are three kinds? (211)=55
E.g.: How many ways are there to park 2 identical cars in 4 parking spots if the cars can stack on top of each other? (35)=10
E.g.: How many different outcomes are there to roll a 6-sided die 10 times if ordering of the outcomes between the trials does not matter (so 123456 is the same as 654321)? (1015)=3003
The way to visualize this problem is to have 10 stars and 6−1=5 bars that separate the stars. So with 0 representing a star and | representing a bar, 0||000|00|0000| would represent one 1, zero 2’s, three 3’s, two 4’s, four 5’s, and zero 6’s. There a total of 10+5=15 “parking spots” and 5 “cars” or 10 “empty spots leftover” depending on how you view it, meaning there are (515)=(1015) possibilities total.
C(n;n1,n2…nk)=(n1n2…nkn)=n1!n2!…nk!n!
This is the number of ways of choosing n items with k distinct groups when repetition is allowed and ordering matters
Also called a multinomial coefficient of a1n1,a2n2,…,aknk for the multinomial expansion of (a1+a2+…+ak)n where n1+n2+…+nk=n
E.g. What is the number of numbers that can be formed using all of the following digits: 111122334 =4!2!2!1!9!=3780 (a+b)n=k=0∑n(kn)an−kbk,{n∈Z:n≥0} (Binomial Theorem)
(kn+1)=(k−1n)(kn),1≤k≤n
Usd to build Pascal’s Triangle
(f:X→Y∧∣X∣>∣Y∣)→∃x1,x2:f(x1)=f(x2)∧x1=x2 (Pigeonhole Principle)
i.e. if there are n pigeons that fly into k pigeon holes where n>k, at least one pigeon hole must have more than one pigeon in it
(f:X→Y,∣X∣=n,∣Y∣=m,k=⌈mn⌉)→∃x1,x2,…,xk:f(x1)=f(x2)=…=f(xk) (Strong Pigeonhole Principle)
i.e. if there are n pigeons that fly into m pigeon holes where k is mn rounded up, at least one pigeon hole must have at least k pigeons in it
Recurrence Relations 🐚
Linear homogenous recurrence relations with constant coefficients
Shortened to LHRRWCC
Linear: The highest degree is one
Homogenous: All terms are of the same degee
Constant coeffcients: Self-explanatory
If some form of tn works, then the linear combination of the solutions to the resulting polynomial are solutions to the recurrence relation.
For order two LHRRWCC (i.e. an=c1an−1+c2an−2), if t2=c1t+c2 has two distinct solutions r1 and r2, the solutions are of the form ∃b,n∈R:br1n+dr2n
Graph G(V,E) describes a graph with vertices and edges
Euler Path
Definition
A path that passes through all edges exactly once
For a Euler path to exist, there must be zero or two nodes that have an odd number of connections. For the two nodes that have an odd number of connections, those are the starting and end points
Euler Cycle
Hamiltonian Cycle
Probability 🎲
Exclusive/Exhaustive Events A1,A2,…,Ak are mutually exclusive events iff Ai∩Aj=∅ when i=j A1,A2,…,Ak are exhaustive events iff A1∪A2∪…∪Ak=Ω, the universal set
Probability Space
A probability space is a triple of (Ω,F,P): Ω is the set of all possible outcomes of some random experiment F is the power set of Ω (all possible events you could try to ask the probability of) P is a probability measure
An event A satisfies: A∈F and A⊆Ω
E.g.: What is the probability of getting at least 1 head if you flipped two coins? Ω={HH,HT,TH,TT} F=P(Ω)
Let A the event that you get at least 1 head A={HH,HT,TH}
Then of course, A∈F and A⊆Ω P(A)=∣Ω∣∣A∣=43
Probability Measure
A probability measure is a function P:F↦[0,1] satisfies
(a) P(A)≥0
(b) P(S)=1
(c) Ai∩Aj=∅,i=j→P(A1∪A2∪…∪Ak)=P(A1)+P(A2)+…+P(Ak) (i.e. mutually exclusive events have cumulative probabilities ) 0≤P(A)=1−P(A′)≤1 P(∅)=0 A⊂B→P(A)≤P(B) P(A\B)=P(A)−P(B)
P(A∪B)=P(A)+P(B)−P(A∩B)
Inclusion-exclusion principle for two events
P(j=1⋃nAj)=k=1∑n(−1)k−1I⊆{1,…,n},∣I∣=k∑P(j∈I⋂Aj)
Generalized inclusion-exclusion principle
Example for 3 events: P(A∪B∪C)=P(A)+P(B)+P(C)−P(A∩B)−P(A∩C)−P(B∩C)+P(A∩B∩C)
A,B are independent events↔P(A∩B)=P(A)P(B)
Generalized form: A1,A2,…,An⊆Ω are mutually independent iff given 1≤k≤n and 1≤j1<j2<…<jn≤n, then P(ℓ=1⋃kAjℓ)=ℓ=1∏kP(Ajℓ)
In other words, each event needs to be pair-wise, tri-wise, …, and n-wise independent
P(A∣B)=P(B)P(A∩B)
This is the conditional probability formula P(A∣B) is read as the probability of A given B P(B)>0 or else B should have been able to occur in the first place (also division by 0) P(A∣B) satisfies the criteria for a probability function
P(A1∩A2∩…∩An)=k∏nP(Ak∣(A1∩…∩Ak−1))
Generalized formula for probability of the union events using conditional probability
(a+b)n=k=0∑n(kn)anbn−k
The famous binomial theorem
Works for n∈N∪{0} and x,y∈R
P(A)=k=1∑nP(Bk)P(A∣Bk)
Called the law of total probability Bk here makes up a partition of S, our sample space
P(A∣B)=P(B)P(B∣A)P(A)
Bayes’ Theorem
Comes from the fact that P(A∩B)=P(B∣A)P(A), which can be plugged into the formula for conditional probability
Distributions 🔔
Random Variables
Definition
For a set S and a probability space (Ω,F,P) X:Ω↦S Ω is the universal set
i.e. a random variable X is a function that assigns outcomes of our random experiments to values from our set S P(X=x):=P({ω∈Ω:X(ω)=x}) P(X∈A):=P({ω∈Ω:X(ω)∈A})
E.g.: Let’s say you play a game where if you flip a coin twice and land at least 1 head, you win 10 dollars, but you otherwise lose 5 dollars. Then we can define a probability space (Ω,F,P): Ω={HH,HT,TH,TT} F=P(Ω) P is a probability function
Then we can define a random variable X such that for ω∈Ω,X(ω)= the number of dollars won for ω, which is X(ω)∈S={−5,10}
FX(x)=P(X≤x)
Cumulative distribution function (CDF) FX:R↦[0,1]
In the two-coin flip example , we may say: FX(10)=P(X=10)=P(winning 10 or less dollars)=P({ω∈Ω:X(ω)=−5})+P({ω∈Ω:X(ω)=10})=P({TT})+P({HH,HT,TH})=43+41=1
FY(x)=FX(x)→X∼Y
If two random variables have the same CDF, then they are said to be identically distributed
Let X be the number of heads that appear when flipping a coin twice and Y be the number of tails that appear when flipping a coin twice
Then X and Y are identically distributed
E[a]=a,a∈R
The expected value of a constant is just the constant
E[ag(X)+bh(X)]=aE[g(X)]+bE[h(X)]
i.e. the expectation of random variables is linear
E[(X−b)r]
This is the rth moment of X about b
An example is that σX2 is the 2nd moment of X about E[X] b=0→ the rth moment of X=E[Xr] σX2=var(X)=E[(X−E[X])2]=E[X2]−E[X]2 σX=STDEV(X) σaX+b2=a2var(X) σaX+b=∣a∣σX
MX(t)=E[etX],t∈R
Called the moment-generating function (MGF) of X E[Xr]=dtrdrMX∣t=0 E[X]=dtdMX∣t=0=dtdlnMX∣t=0 σX2=dt2d2lnMX∣t=0
Let there be g(X,Y) and h(X,Y)
E[ag(X,Y)+bh(X,Y)]=aE[g(X,Y)]+bE[h(X,Y)]
States that expected value demonstrates linearity for bivariate random variables g(x,y)≤h(x,y)∀(x,y)∈S→E[g(X,Y)]≤E[h(X,Y)] E[g(X)]=x∈SX∑g(x)pX(x),E[h(Y)]=y∈SY∑h(y)pY(y) X,Y independent→E[g(X)h(Y)]=E[g(X)]E[h(Y)]
∣E[XY]∣≤E[X2]E[Y2]
Called the Cauchy-Schwarz inequality
cov(X,Y)=E[(X−E[X])(Y−E[Y])]=E[XY]−E[X]E[Y]
Measures how much 2 random variables vary together (is loosely how positively correlated they are)
(1) cov(X,a)=0
(2) cov(X,X)=var(X)
(3) cov(X,Y)=cov(Y,X)
(4) cov(aX,bY)=abcov(X,Y)
(5) cov(X+a,Y+b)=cov(X,Y)
(6) cov(aX+bY,cW+dV)=accov(X,W)+adcov(X,V)+bccov(Y,W)+bdcov(Y,V) X,Y independent→cov(X,Y)=0
ρ(X,Y)=var(X)var(Y)cov(X,Y)
Pearson correlation coefficient
Same idea as r=∑(xi−xˉ)2∑(yi−yˉ)2∑(xi−xˉ)(yi−yˉ), except it also encapsulates pX(x) or fX(x) and pY(y) or fY(y) information into the coefficient ∣ρ(X,Y)∣≤1
g(x)=E[Y∣x],g(X)=E[Y∣X]x` g:SX↦R g(X) is a random variable and is the conditional expectation of Y conditioned on X
h(x)=var(Y∣x),h(X)=var(Y∣X) h:SX↦R h(X) is a random variable and is the conditional expectation of Y conditioned on x
E[E[Y∣X]]=E[Y]
Called the Law of Iterated Expectation
var(Y)=E[var(Y∣X)]+var(E[Y∣X])
Called the Law of Total Variance
Discrete Random Variables
Let X be a discrete random variable
Definition
A random variable is discrete if S from X:Ω↦S is finite (injective onto N)
pX(x)=P(X=x)
Probability mass function (PMF) pX:S↦[0,1] P(X=x):=P({ω∈Ω:X(ω)=x})
In the two-coin flip example , we may say: pX(10)=P(X=10)=P(winning 10 dollars)=P({ω∈Ω:X(ω)=10})=P({HH,HT,TH})=43 P(X∈A)=x∈A∩S∑pX(x) FX(x)=y∈S,y≤x∑pX(y) P(a<X≤b)=FX(b)−FX(a)
μX=E[X]=x∈S∑xpX(x)(if the sum converges)
Called the expected value of the random variable X
X∼Uniform({1,2,…,m})
A random variable is uniformly distrubted if the probability of each x∈S is equal. For the probabilities to be equal and add up to one, and since there are m of x∈S, pX(xi)=m1 pX(x)=m1,x∈{1,2,…,m} FX(x)= MX(t)= E[X]=2m+1 var(X)=
An example of a random variable would be X= the number of pips on a die face after a roll. There are 6 outcomes because m=∣{1,2,3,4,5,6}∣=6, meaning pX(xi)=m1=61
X∼Bernoulli(p) for p∈(0,1)
A Bernoulli random variable takes two values: 0 and 1, which are both specified by the below pX and functions
\begin{cases}
p~~~\qquad\text{if }x=1\\
1-p\quad\text{if }x=0
\end{cases}$$
$$F_X(x) =
\begin{cases}
0 ~~~~~~~\quad \text{if }x < 0\\
1 - p\quad\text{if }0 \le x < 1\\
1 ~~~\qquad \text{if }x \ge 1
\end{cases}$$
$M_X(t) =$
$\mathbb{E}[X]=p$
$\text{var}(X)=p(1-p)$
X∼Binomial(n,p)
A binomial random variable represents the number of successes out of nindependent trials that each have a probability p of success pX(x)=(xn)px(1−p)n−x,x∈{0,1,…,n} FX(x)=
MX(t)=(1−p+pet)n
This occurs because:
(1) MX(t)=x=0∑netxpX(x)
(2) MX(t)=x=0∑netxpx(1−p)n−x
(3) MX(t)=x=0∑n(etp)x(1−p)n−x
(4) By Binomial theorem: MX(t)=(1−p+pet)n E[X]=np var(X)=np(1−p)
A binomial distribution describes the number of successes when running n independent trials where the probability of a success for each trial is p
X∼Geometric(p)
A geometric random variable represents the number of trials until a first success where each trial is independent and has a probability p of success pX(x)=(1−p)x−1p,x∈{1,2,3,…} FX(x)=1−(1−p)x MX(t)=1−(1−p)etpet,t<−ln(1−p) E[X]=p1 var(X)=p21−p
A geometric distribution describes the number of trials until a success where each trial has a probability of p of occuring
X∼Negative Binomial(r,p)
A negative binomial random variable represents the number of trials until the rth success where there each trial has a probability p of success pX(x)=(r−1x−1)pr(1−p)x−r,x∈{r,r+1,…} FX(x)= MX(t)=(1−(1−p)etpet)r,t<−ln(1−p) E[X]=pr var(X)=p2r(1−p)
X∼Poisson(λt)
A Poisson variable represents the number of occurrences of an event in t units of time that has an average occurrence of λ occurrences per unit time
Analogous to binomial distributions in the discrete world
The λ below =λt pX(x)=e−λx!λx,x∈{0,1,2,…} FX(x)= MX(t)=eλ(et−1) E[X]=λ var(X)=λ
Let Y be another discrete random variable pX,Y:S→P(X=x∩Y=y)
pX(x)=y∈SY∑pX,Y(x,y),pY(y)=x∈SX∑pX,Y(x,y) pX(x) and pY(y) are called marginal PMF’s of X and Y respectively pX:SX↦[0,1] pY:SY↦[0,1] x∈SX∑pX(x)=1,y∈SY∑pY(y)=1 X,Y independent↔pX,Y(x,y)=pX(x)pY(y)∀(x,y)∈S E[g(X,Y)]=(x,y)∈S∑g(x,y)pX,Y(x,y)
X∣y→pX∣Y(x∣y)=P(X=x∣Y=y)=pY(y)pX,Y(x,y),x∈SX X∣y is a random variable and contains the above PMF
Continuous Random Variables
Let X be a continuous random variable
fX(x)=P(X=x)
Probability density function (PDF) fX:S↦[0,∞] fX is called the probability density function for X and can be thought of as probability per unit length of the interval. Is akin to pXin the discrete world, except pX(x) is usually not simply 0
FX(x)=∫−∞xfX(t)dt FX:R→[0,1] is still the cumulative distribution function like in the discrete world, where FX(x):=P(X≤x) x→−∞limFX(x)=0,x→∞limFX(x)=1 ∫−∞∞fX(x)dx=1
P(a<X<b)=∫abfX(x)dx
The inequalities don’t have to be strict, they can be ≤ as well because P(X=a)=P(X=b)=0 for continuous random variables E[X]=∫−∞∞xfX(x)dx E[g(X)]=∫−∞∞g(x)fX(x)dx
X∼Cauchy
Called a pathological distribution because most of the important values are undefined fX(x)=π(1+x2)1∀x∈R MX(t)=∅ E[X]=∅ var(X)=∅
X∼Uniform([a,b]) X represents a random point in the interval [a,b]\frac{1}{b - a} \qquad \text{if } x \in (a,b)\\
0 \qquad ~~~~ \text{otherwise}
\end{cases}$$
$$F_X(x) = \begin{cases}
\frac{x - a}{b - a} ~~~ \text{if }x \in (a,b)\\
0 \qquad \text{if }x \le a\\
1 \qquad \text{if }x \ge b
\end{cases}$$
$\mathbb{E}[X] = \frac{b + a}{2}$
$\text{var}(X) = \frac{(b - a)^2}{12}$
X∼Exp(θ) X is the time of the first arrival for a Poisson process with mean number of arrivals per unit time λ>0 and mean waiting time for each arrival θ=λ1>0
Analagous to geometric distributions in the discrete world
\frac{1}{\theta}e^{-\frac{x}{\theta}}, ~ x \ge 0\\
0, ~~~~~~~~ x < 0
\end{cases}$$
$F_X(x) = 1 - e^{-\frac{x}{\theta}}$
$M_X(t) = \frac{1}{1 - \theta t}, ~ t < \frac{1}{\theta}$
$\mathbb{E}[X] = \theta$
$\text{var}(X) = \theta^2$
E.g.: What is the probability that for a store receiving on average 3 customer per minute that a customer comes after just 30 seconds?
$X \sim \text{Exp}(\frac{1}{3}), ~ P(X < \frac{1}{2}) = F_X(\frac{1}{2}) - F_X(0) = 1 - e^\frac{-1/2}{1/3} - (1 - e^{-\frac{0}{1/3}}) = 1 - e^{-3/2} = 1 - 0.223 = 0.777$
X∼Gamma(α,θ) X is the time of the αth arrival of a Poisson process with mean number of arrivals per unit time λ>0 and mean waiting time for each arrival θ=λ1
Analagous to negative binomial distributions in the discrete world Γ(α)=∫0∞xα−1e−xdx Γ(α)=(α−1)!∀α∈Z fX(x)=θαΓ(α)1xα−1e−θx,x>0
Max of fX(x) at x=θ(α−1) MX(t)=(1−θt)α1,t<θ1 E[X]=αθ var(X)=αθ2
X∼χ2(k),k≥1 X is a chi-square distribution with k degrees of freedom
Is a special case of the gamma distribution: χ2(k)∼Gamma(k/2,2) fX(x)=2k/2Γ(k/2)1xk/2−1e−x/2,x>0 FX(x)=Γ(k/2)1γ(2k,2x) γ(s,x)=∫0xts−1e−tdt
X∼N(μ,σ2) μ∈R,σ2>0 fX(x)=σ2π1e−2σ2(x−μ)2,x∈R FX(x)=∫−∞xfX(t)dt MX(t)=eμt+21σ2t2,t∈R E[X]=μ var(X)=σ2
For the standard normal distribution Z=σX−μ∼N(0,1) fZ(z)=2π1e−2z2 FZ(z)=Φ(z)=∫−∞zfZ(t)dt MZ(t)=e21t2 E[Z]=0 var(Z)=σZ=1 Φ(−x)=1−Φ(x)
Table for ϕ(x)
(baby) Central Limit Theorem P(μX−aσX≤X≤μX+bσX)→2π1∫abe−21x2dx as n→∞
P((X,Y)∈A)=∬AfX,Y(x,y)dxdy fX,Y is the joint probability density function ∫−∞∞∫−∞∞fX,Y(x,y)dxdy=1
fY(y)=∫−∞∞fX,Y(x,y)dx
Marginal PDF of Y E[g(X,Y)]=∫−∞∞∫−∞∞g(x,y)fX,Y(x,y)dxdy
fY∣X(y∣x)=fX(x)fX,Y(x,y) ∫−∞∞fY∣X(y∣x)dy=1
Y=u(X)→fY(y)=∣dydu−1(y)∣⋅fX(u−1(y))
The PDF of a random variable that is a function of another variable is given as above u′(x)>0 or u′(x)<0∀x∈S
Several Random Variables
Definitions X1,X2,…,Xn are discrete random variables taking values in sets S1,S2,…,Sn⊆R, S=S1×S2×Sn⊆Rn pX1,X2,…,Xn(x1,x2,…,xn)=P(X1=x1,X2,=x2,…,Xn=xn) pXj(xj)=P(Xj=xj)=x1∈S1∑⋯xj−1∈Sj−1∑xj+1∈Sj+1∑…xn∈Sn∑pX1,X2,…,Xn(x1,x2,…,xn) X1,X2,…,Xn independent ↔pX1,X2,…,Xn(x1,x2,…,xn)=pX1(x1)pX2(x2)⋯pXn(xn) X1,X2,…,Xn independent and identically distriubted (i.i.d) if independent and pXj=pX1,j=2,…n X1,X2,…,Xn are continuous random variables ∃fX1,X2,…,Xn(x1,x2,…,xn):P((X1,X2,…,Xn)∈A)=∫AfX1,X2,…,Xn(x1,x2,…,xn)dx1…dxn fXj(xj)=∫Rn−1fX1,X2,…,Xn(x1,x2,…,xn)dx1…dxj−1dxj−1…dxn fX1,X2,…,Xn(x1,x2,…,xn)=fX1(x1)fX2(x2)⋯fXn(xn)∀(x1,x2,…,xn) X1,X2,…,Xn independent and identically distriubted (i.i.d) if independent and fXj=fX1,j=2,…n E[u(X1,…,Xn)]=(x1,…,xn)∈S∑u(x1,…,xn)pX1,…,Xn(x1,…,xn) E[u(X1,…,Xn)]=∫Rnu(x1,…,xn)fX1,…,Xn)(x1,…,xn)dx1⋯dxn E[au(X1,…,Xn)+bv(X1,…,Xn)]=aE[u(X1,…,Xn)]+bE[v(X1,…,Xn)] u(x1,…,xn)≤v(x1,…,xn)∀(x1,…,xn)→E[u(X1,…,Xn)]≤E[v(X1,…,Xn)] Y=a1X1+…+anXn E[Y]=j=1∑najE[Xj] var(Y)=j=1∑nk=1∑najakcov(Xj,Xk) Xj are independent→var(Y)=j=1∑naj2var(Xj) MY(t)=j=1∏nMXj(ajt)
For i.i.d Xj with mean μ and standard deviation σ Sn=j=1∑nXj Xn=n1j=1∑nXj=n1Sn E[Sn]=j=1∑nE[Xj]=j=1∑nμ=nμ E[Xn]=μ var(Sn)=j=1∑ncov(Xj,Xj)=j=1∑nvar(Xj)=nσ2 var(Xn)=nσ2 (Xn→X in probability as n→∞)↔P(∣X−Xn∣≥ε)→0 as n→∞ for any ε
Xn=n1j=1∑nXj→μ as n→∞
Weak law of large numbers
The sample average random variable approaches in probability the expected value of each sample as we increase the number of samples
P(∣X−μ∣≥λ)≤λ2σ2
Chebyshev’s Inequality λ>0 MX(t)=MY(t)→X,Y i.i.d. MSn(t)=[M(t)]n MXn(t)=[M(nt)]n (Xn→X in distribution as n→∞)↔FXn(x)→FX(x) as n→∞ (Xn→X in distribution as n↔MXn(t)→MX(t) as n→∞,h>0∩t∈(−h,h) (Xj∼N(μj,σj2),Y=j=1∑najXj)→Y=N(j=1∑nμjaj,j=1∑nσj2aj2) var(Xn)=E[(Xn−μ)2]=nσ2
Central Limit Theorem: σ/nXn−μ→N(0,1) in distribution as n→∞
States that the sample average random variable approaches a normal distribution given a large enough sample size
3B1B Video: https://www.youtube.com/watch?v=zeJD6dqJ5lo
“Sample size” here is not to be confused with the number of simulations taken when we are trying to sus out a normal distribution curve. So at ~7:04 in the video, # Sums refers to the number of simulations, but the sample size here is only 10.
Results: Xn≈μ+nσZ Xn≈N(μ,nσ2) Sn=j=1∑nXj≈N(nμ,nσ2)
max(a,b)=2a+b+(2a+b)2−ab=2a+b+∣a−b∣min(a,b)=2a+b−(2a+b)2−ab=2a+b−∣a−b∣
This property can be used to solve for a or b if all we know are their sum and product, which is useful for factoring.
E.g. to factor x2+52x+672, we do: a=252+(252)2−672=26+676−672=26+2=28b=52−a=52−28=24
Factoring Algorithm for ax2+bx+c
E.g. for 2x2+3x+1
(1) 21((2x)2+3(2x)+2)
(2) 21(2x+1)(2x+2)
(3) (2x+1)(x+1)
Simplified: ax2+bx+c=a1(u2+bu+ua),u=ax
Ar⋅Bs=⟨AB⟩∣r−s∣
The dot product decreases the grade of the multivector from r and s to r−s
e.g. The grade of a⋅b for a,b∈V=1−1=0
Ar∧Bs=⟨AB⟩r+s
The wedge product increases the grade of the multivector from r and s to r+s
e.g. The grade of a∧b for a,b∈V=1+1=2, which is the grade of a bivector
anf=ani+d(nf−ni);an=a1+d(n−1)
This is the formula for the next element in an arithmetic sequence anf is the nfth element of the sequence while ani is the nith element of the sequence. d is the difference between each element in the sequence. ai by default starts at 0.
An example would be {2, 5, 8, 11, 14, …}, where d=3 and a1=2. Then the 6th element of this sequence would be anf=2+3(6−1)=17
S=2nf−ni+1(anf+ani),nf−ni=danf−ani S is the sum of elements of an arithmetic sequence with anf being the last element and ani being the first element in the series
The visual intuition is to imagine a grid staircase with the first step being ani high and the last step being anf high with n number of total steps. Duplicating the staircase and aligning them together to form a rectangle allows you to count the total number of squares:
Desmostartion of the formula: https://www.desmos.com/calculator/1rofjxm3op
n=danf−ani+1 n is the number of elements in the arithmetic series
Geometric
anf=anir(nf−ni);an=a1rn−1
This is the formula for the next element in a geometric sequence anf is the nfth element of the sequence while ani is the nith element of the sequence. r is the ratio between the (n+1)th element and the nth element the sequence. a0 by default starts at 1.
An example would be {2, 6, 18, 54, 162, …}, where r=3 and a1=2. Then the 6th element of this sequence would be anf=2⋅36−1=486
S=r−1ani(rnf−ni+1−1),r=1;S∞=1−rani,∣r∣<1 S is the sum of elements in a geometric series from an0 to anf
Derivation: Snf=an0+an0r+…+an0rnf−n0+1 Snfr=Snf−an0+an0rnf−n0 Snf(r−1)=an0rnf−n0−an0 Snf=r−1an0(rnf−n0−1),r=1
Desmostration: https://www.desmos.com/calculator/quepjw8doa
n=ln(r)ln(anianf)+1 n is the number of elements in the geometric series n=1∑Nn=2n(n+1) n=1∑Nn2=6n(n+1)(2n+1) n=1∑Nn3=(2n(n+1))2
f(x)=f(a)+f′(a)1!(x−a)+f′′(a)2!(x−a)2+f′′′(a)3!(x−a)3+…=n=0∑∞f(n)(a)n!(x−a)n
Taylor Series
Maclaurin series is when a=0 https://www.youtube.com/watch?v=LDBnS4c7YbA ex=k=0∑∞k!xk=1+x+2!x2+3!x3+4!x4+…
Mass and Probability M=mass,δ(x,y)=mass density, M=∬Dδ(x,y)dA My=∬Dxδ(x,y)dA (y-moment) Mx=∬Dyδ(x,y)dA (x-moment) (xˉ,yˉ)=(MMy,MMx) Ix=∬Dx2δ(x,y)dA (moment of inertia about x-axis) Iy=∬Dy2δ(x,y)dA (moment of inertia about y-axis) IO=∬D(x2+y2)δ(x,y)dA M=mass,δ(x,y,z)=mass density, M=∭Wδ(x,y,z)dV Myz=∭Wxδ(x,y,z)dV (yz-moment) Mxz=∭Wyδ(x,y,z)dV (xz-moment) Mxy=∭Wzδ(x,y,z)dV (xy-moment) (xˉ,yˉ,zˉ)=(MMyz,MMxz,MMxy) Ix=∭W(y2+z2)δ(x,y,z)dV (moment of inertia about x-axis) Iy=∭W(x2+z2)δ(x,y,z)dV (moment of inertia about y-axis) Iz=∭W(x2+y2)δ(x,y,z)dV (moment of inertia about z-axis) IO=∭W(x2+y2+z2)δ(x,y,z)dV P=total probability,pX(x)=probability density function for continuous random variable X P=∫−∞∞pX(x)dx=1
P[a<X≤b]=∫abpX(x)dx
This is the probability that a<X≤b
E[f(X)]=∫−∞∞f(x)pX(x)dx
This is the expected value of f(X)
Only applies if f:R→R P=total probability,pX,Y(x,y)=probability density function for continuous random variables X and Y P=∬−∞∞pX,Y(x,y)dxdy=1
P[(X,Y)∈D]=∬DpX,Y(x,y)dA
This is the probability that (X,Y)∈D
E[f(X,Y)]=∬−∞∞f(x,y)pX,Y(x,y)dxdy
This is the expected value of f(X,Y)
Only applies if f:R2→R pX(x)=∫−∞∞pX,Y(x,y)dy pY(y)=∫−∞∞pX,Y(x,y)dx
Operations on Vector Fields ∇=⟨∂x∂,∂y∂,∂z∂⟩(Del Operator)
Δ=∇2=∇⋅∇=∂x2∂2+∂y2∂2+∂z2∂2 (Laplace operator)
The Laplacian of a function is akin to the second derivative of a function in 1D calculus. However, there is now more than one input. So Δf=∇⋅(∇f), which is the divergence of the gradient of f. So more positive Δf corresponds to minima while more negative Δf corresponds to maxima.
3Blue1Brown video that explains the intuition of the Laplacian: https://www.youtube.com/watch?v=EW08rD-GFh0
∇f(x,y,z)=∂x∂f∂y∂f∂z∂f=F, which is a vector field (not the same as a vector)
The gradient of a scalar field evaluated at a point gives a vector that points toward greatest increase of the scalar field
\frac{\partial f}{\partial y}+
\frac{\partial f}{\partial z}$$
$$\text{curl}(\vec F)=\nabla\times\vec F=\begin{vmatrix}
\hat i&\hat j&\hat k\\
\frac{\partial f}{\partial x}&\frac{\partial f}{\partial y}&\frac{\partial f}{\partial z}\\
F_1&F_2&F_3
\end{vmatrix}\\=(\frac{\partial F_3}{\partial y}-\frac{\partial F_2}{\partial z})\hat i+(\frac{\partial F_1}{\partial z}-\frac{\partial F_3}{\partial x})\hat j+(\frac{\partial F_2}{\partial x}-\frac{\partial F_1}{\partial y})\hat k$$
- $\vec F\text{ is conservative in 2D }\rightarrow \frac{\partial F_1}{\partial y}=\frac{\partial F_2}{\partial x}$
This statement is biconditional if $F$ is on a simply connected domain (no holes)
- $\vec F\text{ is conservative in 3D}\rightarrow\nabla\times F=0$
This statement is biconditional if $\vec F$ is on a simply connected domain (no holes that prevent any choice of a closed curve from being stretched and deformed to a point without leaving the domain)
Scalar, Vector, and Flux Integrals
Scalar line integral:∫Cf(x,y,z)ds=∫t0tff(r(t))∣∣r′(t)∣∣dt
If f represents mass density and ds represents some infinitesimal length of the curve, the scalar line integral would represent the mass of the curve segment
Might involve finding a parameterization of f in terms of t
Ex: What is the mass of a wire represented by the curve y=x3,{0≤x≤1} (cm) if its mass density is represented by f(x,y)=1+9xy (g/cm) r(t)=⟨t,t3⟩ M=∫C1+9xyds=∫011+9t4∣∣⟨1,3t2⟩∣∣dt=514g
Vector line integral:∫CF⋅tds=∫CF⋅dr=∫C(F1dx+F2dy+F3dz)=∫t0tfF(r(t))⋅r′(t)dt t represents the unit tangent vector of the path
If F represents a vector field for some force, then the vector line integral represents the work done by the force dx=x′(t)dt,dy=y′(t)dt,dz=z′(t)dt, where r′(t)=⟨x′(t),y′(t),z′(t)⟩ ∫C∇f⋅dr=f(Sf)−f(S0) (Fundamental Theorem of Vector Integrals)
2D Flux:∫CF⋅n^ds=∫t0tfF(r(t))⋅N(t)dt n^ is obtained by rotating t^ by 2π clockwise t can be simplified to just r′(t)
Vortex Field F(x,y)=⟨−x2+y2y,x2+y2x⟩ if (x,y)=(0,0)
Satisfies the cross-partials condition
Not conservative on paths encircling the origin G(u,v)=(x(u,v),y(u,v),z(u,v))
Tu(P)=∂u∂G(u,v),Tv(P)=∂v∂G(u,v) Tu points toward increasing u, Tv points toward increasing v, N(P)=Tu(P)×Tv(P) Nsphere(θ,ϕ)=⟨−cosθsin2ϕ,−sinθsin2ϕ,−cosϕsinϕ⟩ (⟨x,y,z⟩−G(u,v))⋅N=0 (Equation for tangent plane)
Area(S)=∫v0vf∫u0uf∣∣N(u,v)∣∣dudv N(u,v) is the same as ∣∣Tu×Tv∣∣, which represents the area of a small paralleleogram on the surface S. This representation happens because ∣∣Tu×Tv∣∣ is ∣∣Tu∣∣∣∣Tv∣∣sin(θ), where θ is the angle between the two vectors that make up the parallelogram. Tu for example, is ∂u∂GΔu=G(u+Δu,v)−G(u,v)=u
Scalar surface integral:∬Sf(x,y,z)dS=∫v0vf∫u0uff(G(u,v))∣∣N(u,v)∣∣dudv ∬CdS=Area(s)
Is akin to scalar line integral in that it is the mass of the surface if f is mass density
Vector surface integral/3D Flux: ∬SF⋅n^dS=∬SF⋅dS=∫v0vf∫u0ufF(G(u,v))⋅N(u,v)dudv
Physically measures the amount of “stuff” (say electric field or water current) that passes through a surface
Is akin to the vector line integral and 2D flux, but for surfaces.
Heuristically, dS=Ndudv,dS=∣∣N∣∣dudv
Green’s Theorem:∬Dcurlz(F)dA=∮CF⋅dr curlz(F(x,y))=∂x∂F2−∂y∂F1. It is the z-component of the curl of F C is the boundary of D oriented positively, or “going counter clockwise”
The intuition is that C is a closed curve going counter clockwise, and curlz(F) represents the amount of counter clockwise rotation around a point. The total “alignment” between the vector field F and the counter clockwise boundary C can be obtained by summing up the curls of all the points within our boundary. This idea is explained in this video: https://www.youtube.com/watch?v=8SwKD5_VL5o
Green’s Theorem (flux form): ∬Ddiv(F)dA=∮∂DF⋅n^ds=n=1∑N∮CnF⋅n^ds
The intuition is that the amount of stuff (e.g. fluid) represented by F leaving a point is represented by div(F). The total amount of fluid leaving an entire domain would thus be the sum of divergences of all points within the domain and also the flux integral about the path of the domain’s boundary.
Mean value property: f(x,y) is harmonic→∀(x,y)∈D,f(x,y)=perimeter1∮∂Df(x,y)ds
A function is harmonic iff Δf=0
This essentially states that if a function is harmonic, the value of the function within any domain is equal to the average value of the domain’s perimeter.
∬Df div(F)dA=∮∂DfF⋅n^ds−∬D∇f⋅FdA
Analagous to integration by parts in 1D, but for 2D vectors ∫abf(x)g′(x)dx=[f(x)g(x)]x=ax=b−∫abf′(x)g(x)dx
∬DfΔgdA=∮∂Df∇g⋅n^ds−∬D∇f⋅∇gdA
Also called Green’s Formula
Another form of integration by parts in 2D.
Proof: Take F=∇g and apply integration by parts in 2D
Stoke’s Theorem: ∬Scurl(F)⋅dS=∮∂SF⋅dr=∫u0uf∫v0vfcurl(F(G(u,v)))⋅N(u,v)dvdu
Is like Green’s Theorem, but for 2D surfaces in 3D space. ∂S is the positively oriented boundary of S, meaning you apply the right hand rule such that if the normal points out of the surface, the thumb points toward the surface while the curled fingers dictate the positive orientation of ∂S. Another way of viewing this is to imaging you are on the boundary with your body pointing upright in the same direction as the normal. If you follow the direction of ∂S, then your right foot should be closer to off the edge of the boundary
∮Scurl(F)⋅dS=0
States that the curl of a function over a closed surface is 0. F has a vector potential ↔F=curl(A)→div(A)=0 curl(∇f)=0,div(curl F))=0
Divergence Theorem: ∭Wdiv(F)dV=∮SF⋅dS
Also called Gauss’s Theorem
Like 3D version of Stoke’s Theorem, where S is a positively oriented boundary of W. That essentially means the normal of S points outward.
∭Wfdiv(F)dV=∮∂WfF⋅dS−∭W∇f⋅FdV
Integration by parts in 3D
Proof: apply product rule to ∭Wdiv(fF)dV=∮∂WfF⋅dS
a_{11}&\cdots&a_{1m}\\
\vdots&\ddots&\vdots\\
a_{n1}&\cdots&a_{nm}\\
\end{bmatrix}$$
A system of linear equation’s coefficients represented by a coeffecient matrix would have $n$ equations and $m$ variables
- $\text{Elementary Row Operations}$
(1) Scale a row by a nonzero factor
(2) Shuffle rows
(3) Add nonzero multiple of one row to another row
These are the only operations allowed when solving systems in matrix form
- $\text{Reduced row-echelon form}$
(i) The first nonzero entry must be a 1 in a nonzero row, called the leading 1/pivot
(ii) If a column has a pivot, all other entries in the column must be 0
(iii) All rows with pivots must contain a pivot above and to the left
E.g.:
$$\begin{bmatrix}
0&1&0&3&5\\
0&0&1&3&1\\
0&0&0&0&0
\end{bmatrix}$$
Finding $\text{rref}(A)$ of your augmented matrix $A$ of your system of equations results in a solution to the system where the above matrix would represent:
$$\begin{cases}
x_2+3x_4=5\\
x_3+3x_4=1
\end{cases}$$
- $\text{Consistency}$
A system is consistent iff there is 1 or infinite solutions to it; otherwise, the system is inconsistent
A linear system is inconsistent iff $\text{rref(A)}$ has the row:
$$\begin{bmatrix}
0 & 0&\cdots&1
\end{bmatrix}$$
- $\text{Rank}$
$\text{rank(A) = number of leading 1's in rref(A)}$
$\text{(1) rank(A)}\le n$
$\text{(2) rank(A)}\le m$
$\text{(3) System is inconsistent}\rightarrow\text{rank(A)}\lt n$
If $A\vec x$ has only one solution, then
$\text{rank}(A)=m$
- $A=n\times p,\space B=p\times m,\space AB=n\times m$
$$AB=\begin{bmatrix}
A\vec v_1&A\vec v_2&\cdots&A\vec v_m
\end{bmatrix},\space B=\begin{bmatrix}
\vec v_1&\vec v_2&\cdots&\vec v_m
\end{bmatrix},\space A\vec v_i=v_{i_1}A_{j1}+v_{i_2}A_{j2}+\ldots + v_{i_m}A_{jm}$$
$T:\mathbb{R}^m\rightarrow\mathbb{R}^n,\space A=n\times m$
- $\text{Properties of Matrix Multiplication}$
Not commuative $AB \ne BA$
Associative $A(BC) = (AB)C$
Has identity, $AI_p=I_nA=A$
Distributive for matrix and scalar multiplication $A(B+C)=AB+AC \,\forall A\in\mathbb{M},\mathbb{R}$
Transformations
Linear Transformations (1)A(x+y)=Ax+Ay (2)A(kx)=kA(x)
Linear transformations necessarily follow (1) and (2)
Another definition is that a function T(x) is a linear transformation if there exists an n×m matrix A such that T(x)=Ax for all x in Rm
scale(x)=[k00k]x k is the scaling factor
\begin{bmatrix}
u_1^2&u_1u_2\
u_1u_2&u_2^2
\end{bmatrix}\vec x
=\begin{bmatrix}
u_1^2&u_1u_2&u_1u_3\
u_1u_2&u_2^2&u_2u_3\
u_1u_3&u_2u_3&u_3^2
\end{bmatrix}\vec x This is the projection of $\vec x$ onto $\hat u$, and it is akin to casting a light perpendicular to $\hat u$ so that $\vec x$ casts a shadow on $\hat u$
Another useful form is\frac{1}{w_1^2+w_2^2}\begin{bmatrix}
w_1^2&w_1w_2\
w_1w_2&w_2^2
\end{bmatrix}\vec x$$ w is a non-normalized vector
a&b\
b&-a
\end{bmatrix}\vec x,\space a^2+b^2=1 This is the reflection of $\vec x$ about a vector $\vec u$
Another useful form is $(2A-I_2)\vec x$, where $A$ is the coefficient matrix from $\text{proj}_{\vec u}\vec x$
The way to think about reflections are that\begin{bmatrix}
a\
b
\end{bmatrix}isthereflectionof\begin{bmatrix}
1\
0
\end{bmatrix},and\begin{bmatrix}
-b\
a
\end{bmatrix}isthereflectionof\begin{bmatrix}
0\
1
\end{bmatrix}$$
\cos\theta&-\sin\theta\
\sin\theta&\cos\theta
\end{bmatrix}\vec x=
\begin{bmatrix}
a&-b\
b&a
\end{bmatrix}\vec x,\space a^2+b^2=r^2
$\theta$ is the rotation from the x-axis toward the y-axis
$r$ is the scaling factor of the vector, with 1 being no scaling
One can imagine the transformation as scaling and rotating the vector\begin{bmatrix}
1\
0
\end{bmatrix}to\begin{bmatrix}
a\
b
\end{bmatrix}$$, transforming the curve/coordinate plane together as well
\begin{bmatrix}
1&0\
k&1
\end{bmatrix}\vec x$$
A vertical shear transformation would keep rectangle’s vertical sides vertical k is the shearing factor
\begin{bmatrix}
1&k\
0&1
\end{bmatrix}\vec x$$
A hortizontal shear transformation would keep a rectangle’s horizontal sides horizontal
Invertibility
Rules
For A=n×n, either all hold true or none do:
∃A−1
∃!x∈Rn:Ax=b∀b∈Rn
rref(A)=In
rank(A)=n
im(A)=Rn
ker(A)={0}
Column vectors of A form a basis of Rn
Column vectors of A span Rn
Column vectors of A are linearly independent
detA=0
0 not an eigenvalue of A
To find A−1, you have to isolate x1,…,xm in terms of y1,…,ym
i.e. compute rref[A∣In]. If it’s [In∣B], then A−1=B ∃A−1→rref[A∣In]=[In∣A−1] (AB)−1=B−1A−1
im(T)={T(x):x∈Rm}=span(v1,…,vm) T:Rm→Rn
Same as the range of a function im(T)⊆Rn and im(T) is a subspace of Rn
Also known as the column space of A, where A is a matrix of column vectors representing T
ker(T)={x∈Rm:T(x)=0} T:Rm→Rn
Can be thought of as the zeroes of a polynomial function im(T)⊆Rn and ker(T) is a subspace of Rm
Also called the null space of A, where A is a matrix of column vectors representing T
Bases and Dimensions
Basis
The basis of a span is the smallest linearly independent list of vectors that constitute the span
span(v1,…,vm)={c1v1+…+cmvm:c1,…,cm are scalars}
The span of some vectors is the set of all possible linear combinations of the vectors
V∈Rn→dim(V)=m≤n dim(V) is the number of vectors in the basis of V
There are at most m linearly independent vectors in V
We need at least m vectors to span V
If m vectors in V are linearly independent, they form a basis of V
If m vectors in V span V, then they form a basis of V
dim(im(A))=rank(A)
The basis of the image of a matrix is the set of the columns containing the leading variables. Can be found by taking the rref(A) and then looking at the corresponding columns with a leading 1 (can’t look at the matrix from rref(A) specifically, have to look at A)
dim(ker(A))=m−rank(A)=nullity(A)
The basis of the kernel of matrix is the set of the columns containing the free variables. Can be found by taking the rref(A), finding x in terms of the free variables and then taking the vectors that arise as a result of isolating the free variables as the basis vectors
nullity(A)+rank(A)=m
Called the Rank-nullity theorem
Is a portion of the fundamental theorem of linear algebra
Coordinates
\begin{bmatrix}
c_1\
c_2\
\vdots\
c_m
\end{bmatrix}
\rightarrow
\vec x = c_1 \vec v_1 + c_2 \vec v_2 + \ldots + c_m \vec v_m$$ B is some basis (v1,v2,…,vm) that spans V in Rn [x]B is the B-coordinate vector of some vector x in V c1,c2,…,cm are the B-coordinates of x
(1) [x+y]B=[x]B+[y]B∀x,y∈V
(2) [kx]B=k[x]B,∀x∈V
AS=SB,B=S−1AS,A=SBS−1 A is the standard matrix for T(x), meaning its columns represent how ei are tranformed. It converts x to T(x) B is the B-matrix for T(x), meaning its columns represent how vi in V∈Rn are transformed. It converts [x]B to [T(x)]B\begin{bmatrix}
\vec v_1 \ldots \vec v_n
\end{bmatrix}$$. It converts $[\vec x]_\mathfrak{B}$ to $\vec x$ and $[T(\vec x)]_\mathfrak{B}$ to $T(\vec x)$

If the any of the statements are true and $S$ is invertible, $A$ and $B$ are considered to be similar
Similarity is an equivalence relation (reflexive, symmetric, and transitive)
Diagonal B
Orthogonality
Orthonormal Vectors
Vectors are orthonormal if their magnitude is 1 (normalized) and they are all perpendicular/orthogonal to each other
Orthonormal vectors are linearly independent and n of them can form a basis of Rn
projV(x)=x∥=i=1∑m(ui⋅x)ui
The projection of x onto a subspace V is equal to the sum of how much x aligns with each basis vector of V
In the case where m=dim(Rn), the projection is just equal to x
V⊥={x∈Rn:v⋅x=0∀v∈R}
The orthogonal complement of V is all x that are perpendicular to V
It is also the kernel of projV(x)
(1) V⊥∈R
(2) V∩V⊥=0
(3) dim(V)+dim(V⊥)=n
(4) (V⊥)⊥=V
∥x+y∥2=∥x∥2+∥y∥2↔x⋅y=0
Pythagorean theorem
The square of the magnitude of the the sum of two vectors is equal to the sum of the squares of the magnitude of the vectors iff the vectors are perpendicular to each other
∣x⋅y∣=∥x∥∥y∥cos(θ)
A consequence is the Cauchy-Schwarz inequality: ∣x⋅y∣≤∥x∥∥y∥
r=cos(θ)=∥x∥∥y∥x⋅y
The correlation coefficient r between two variables is the depends on the angle θ between the variables’ deviation vectors x and y
The deviation vectors are how far each point is from the average for the given variable r=(∑(xi−xˉ)2∑(yi−yˉ)2)∑(xi−xˉ)(yi−yˉ) is the usual formula, but they are the same thing
Gram-Schmidt Process
Is an algorithm to convert from a basis to an orthonormal basis (one where each basis vector is orthogonal to each other and has magnitude 1
For each vector vi after the first vector v1, the orthogonal version of vi⊥ relative to span(v1,…,vi−1) is equal to vi−vi∥, where vi∥=projspan(v1,…,vi−1)vi=j=1∑i−1(uj⋅vi)uj
The vectors’ orthogonal vectors then need to be normalized by dividing by their magnitude
M=QR M represents the B-matrix, where B is just some basis Q represents the U-matrix, where U is the orthonormal version of B R represents the change of basis matrix from B to U rij=0 for i>j\begin{bmatrix}
\Vert\vec v_1\Vert & \vec u_1\cdot\vec v_2 & \cdots & \vec u_1\cdot\vec v_1\\
0 & \Vert\vec v_2^\perp\Vert & \ddots & \vdots\\
\vdots & \ddots & \ddots & \vec u_{n-1}\cdot \vec v_n\\
0 & \cdots & 0 & \Vert\vec v_n^\perp\Vert
\end{bmatrix}$$
Orthogonal Transformations/Matrices
An orthogonal transformation preserves absolute lengths and absolute angles
(1) AB is orthogonal if A and B are each orthogonal and they are n×n
(2) A−1 is orthogonal if A is orthogonal
(a) T mapping Rn→Rn iff T(e1),…,T(en) form an orthonormal basis of Rn
(b) n×n matrix is orthogonal iff its columns are orthonormal
AT
The ijth entry of the transpose AT of A is the jith entry of A Symmetric: AT=A Skew-symmetric: AT=−A
(a) (A+B)T=AT+BT
(b) (kA)T=kAT
(c) (AB)T=BTAT
(d) rank(AT)=rank(A)
(e) (AT)−1=(A−1)T
v⋅w=vTw
Note: not the same as vwT, which would be an n×n matrix instead A is an n×n orthogonal matrix↔A−1=AT
\vec u_1 & \ldots \vec u_m
\end{bmatrix}$$
Determinants
Determinant
When detA=0, A is invertible
The determinant of upper or lower triangular matrices or diagonal matrices is just the product of the main diagonal elements
Swapping rows or columns results in negating the determinant
detA=ad−bc=Areaparallelogram
This is the parallelogram formed by column vectors from A
\vec v_1\
\vdots\
\vec x\
\vdots\
\vec v_n
\end{bmatrix},\space L_2(\vec x)=\det\begin{bmatrix}
\vec v_1 & \cdots & \vec x & \cdots & \vec v_n
\end{bmatrix}$$
The above statements says that the functions L1 and L2 from Rn→R are linear, meaning that the determinant of a matrix can be written as the sum of the determinants of two matrices whose one of rows or columns add together to form the original matrix
The matrix also doesn’t have to be square
E.g.: \vec u & \vec v & \vec w
\end{bmatrix}\text{ is invertible}$$
- $\det A=\sum(\text{sgn }P)(\text{prod }P)$
$P$ is the pattern or permutation of a matrix. It is defined as a choice of elements in a matrix that do not overlap in rows or columns (so kind of like in Sudoku)
There are $n!$ patterns, so $n!$ terms need to be summed up
$\text{sgn }P=(-1)^p$, where $p$ is the number of inversions of a pattern to reach diagonal or max number of elements that are above and right to one of them
$\text{prod }P$ is simply the product of the pattern

- $\det\begin{bmatrix}
A & B\\
0 & C
\end{bmatrix}=\det\begin{bmatrix}
A & 0\\
B & C
\end{bmatrix}=(\det A)(\det C)$
The determinant of a triangular block matrix is the same as the determinant of the matrix containing the determinant of each block, which would just be
$$\det\begin{bmatrix}
\det A & \det B\\
0 & \det C
\end{bmatrix}$$
for example.
The determinant of a block matrix is not generally just the determinant of a matrix containing the determinant of each block, as the block matrix has to be triangular
$\det A^T=\det A$
$\det A^{-1}=\frac{1}{\det A}$
$\text{Deterimants of Matrices After Row Operations}$
$\text{(1) Multiplying row by } k\text{: }\det B=k\det A$
$\text{(2) Swapping rows: }\det B=-\det A$
$\text{(3) Adding a scalar multiple of a row to another: }\det B=\det A$
$\text{For }n\times n\text{ matrices } A\text{ and }B\text{:}$
$\det(AB)=(\det A)(\det B)$
$\det(A^m)=(\det A)^m$
$A\text{ similar to }B\text{ by an invertible matrix }S\text{: }\det A=\det B$

- $\det T=\det B$
$B$ is the $\mathfrak{B}$-matrix of $T$, which is a linear transformation represented by some matrix $A$
Eigenvalues and Eigenvectors
Av=λv v is called an eigenvector of A or T λ is called an eigenvalue of A or T
B=(vi) is an eigenbasis↔Avi=λivi
An eigenbasis is a basis with basis vectors that are all eigenvectors
A is diagonalizable↔∃S:S−1AS=B is diagonal
A square transformation matrix is diagonalizable if there exists a basis in Rn such that [T]B is diagonal
In other words, A is diagonalizable if it’s similar to a diagonal matrix
Diagonalizing matrices makes certain calculations easier (e.g. A35 is more easily computed as (SBS−1)35=SB35S−1, where B is a diagonal matrix)
B=(vi) is an eigenbasis for A↔S−1AS=B diagonalizes A\vec v_1 & \vec v_2 & \cdots & \vec v_n
\end{bmatrix},\space B=\begin{bmatrix}
\lambda_1 & 0 & \cdots & 0\\
0 & \lambda_2 & \ddots & \vdots\\
\vdots & \ddots & \ddots & 0\\
0 & \cdots & 0 & \lambda_n
\end{bmatrix}$$
$AS=SB,\space A=SBS^{-1}$
The vectors in $S$ must be linearly independent since $S$ is invertible
Conceptually, the vectors of $S$ must form a basis spanning $\mathbb{R}^n$ because if $A=SBS^{-1}$, we are going from the standard basis to some eigenbasis, transforming it by $B$, and then returning back to the standard basis. Thus, there must be $n$ linearly independent vectors of $S$
for the eigenbasis to be traversable from $A$ and back
λi for orthogonal matrix=1 or −1e^{\lambda_1 t} & \cdots & 0\\
\vdots & \ddots & \vdots\\
0 & \cdots & e^{\lambda_n t}
\end{bmatrix}S^{-1}$$
- $\vec x(t+1)=A\vec x(t),\space \vec x(0)=\vec x_0\rightarrow\vec x(t)=\displaystyle\sum_{i=1}^n c_i\lambda_i^t\vec x_0$
For a discrete dynamical system (like differential equations with discrete time inputs), the closed solution is given as the above summation
The $n\times n$ matrix, $A$, has eigenbasis $\mathfrak{B}=(\vec v_i)$ and eigenvalues $\lambda_i$
$\vec x_0=\displaystyle\sum_{i=1}^{n}c_iv_i$
In other words, each $c_i$ is the coefficient for the linear combination of $\vec x_0$ in terms of the eigenvectors $\vec v_i$
- $\vec x(-i)=(A^i)^{-1}\vec x_0$
To look at a system before $t=0$, use this formula
- $\text{Phase Portraits}$
These are curves obtained by connecting the dots of discrete dynamical systems
- $\exists\lambda\text{ of }A\leftrightarrow\det(A-\lambda I_n)=0$
Called the characterisitc or secular equation of $A$
$\text{Eigenvalues of triangular matrices are the diagonal entries}$
- $f_A(\lambda)=\det(A-\lambda I_n)=(-\lambda)^n+\text{tr }A(-\lambda)^{n-1}+\ldots+\det A$
Called the characteristic polynomial of $A$
$\text{tr }A$ is the sum of the diagonal entries of $A$
An example for a $2\times 2$ $A$:
$f_A(\lambda)=\det(A-\lambda I_2)=\lambda^2-(\text{tr }A)\lambda+\det A$
- $f_A(\lambda)=(\lambda_0-\lambda)^kg(\lambda)\land g(\lambda_0)\ne0\leftrightarrow\text{almu}(\lambda_0)=k$
If the characteristic polynomial of a matrix $A$ can be written as the product of a linear factor of $\lambda_0$ with algebraic multiplicity $k$ times some function $g(\lambda)$ that does not have $\lambda_0$ as a solution, then the eigenvalue $\lambda_0$ has an algebraic multiplicty of $k$
- $\#\lambda_\text{distinct}\le n,\space n\text{ is odd}\rightarrow\#\lambda_\text{real}\ge1$
The number of distinct eigenvalues of an $n\times n$ matrix $A$ is at most $n$
If $n$ is odd, then there will be at least one real eigenvalue of $A$
$\displaystyle\prod_\text{i=1}^{\#\lambda}\lambda_i=\det A$
$\displaystyle\sum_{i=1}^{\#\lambda}\lambda_i=\text{tr }A$
- $E_\lambda=\text{ker}(A-\lambda I_n)=\{\vec v\in\mathbb{R}^n:A\vec v=\lambda\vec v\}$
Called the eigenspace associated with $\lambda$
All eigenvectors of $A$ with eigenvalue $\lambda$ are contained in $E_\lambda$
- $\text{gemu}(\lambda)=\text{dim}(E_\lambda)=\text{nullity}(A-\lambda I_n)=n-\text{rank}(A-\lambda I_n)$
The geometric multiplicity of an eigenvalue $\lambda$ is the dimension of its corresponding eigenspace $E_\lambda$
To find the dimension, use rank-nullity theorem
- $\text{Eigenbases and Geometric Multiplicities}$
Given $s=\displaystyle\sum_{i=0}^N\text{gemu}(\lambda_i)$ for an $n\times n$ matrix $A$:
(1) Concatenating the basis vectors of each eigenspace will give a new basis of dimension $s$, implying $s\le n$
(2) $A$ is diagonalizable iff $s=n$
Here, $N$ is the number of distinct eigenvalues
$n\text{ distinct eigenvalues}\rightarrow A\text{ is diagonalizable}$
$AS=SB\rightarrow$
(a) $f_A(\lambda)=f_B(\lambda)$
(b) $\text{rank }A=\text{rank }B,\space \text{nullity }A=\text{nullity }B$
(c) $A \text{ and }B\text{ have the same eigenvalues, almu, and gemu}$
(d) $\det A=\det B,\space \text{tr }A=\text{tr }B$
$\text{gemu}(\lambda)\le\text{almu}(\lambda)$
Complex Linear Algebra
Operations on complex numbers z,z1,z2,w∈C:
(1) Commutativity of addition: z+w=w+z
(2) Commutativity of multiplication: zw=wz
(3) Associativity of addition: (z1+z2)+z3=z1+(z2+z3)
(4) Associativity of multiplication: (z1z2)z3=z1(z2z3)
(5) Distributivity: z(w1+w2)=zw1+zw2
z=a+ib,zˉ=a−ib
The conjugate of a complex number negates the imaginary component of the number z=a+ib,z−1=z1=(a2+b2a)+i(a2+b2−b)=a2+b2a−ib
A with vectors ∈C has n distinct complex eigenvalues if they are counted with their algebraic multiplicities
“Counted with their algebraic multiplicities” means you count three times for example if λi has degree three
Orthonormal Eigenbasis A has an orthonormal eigenbasis S if S−1=ST (S is orthonormal) and A=SBS−1 where B is a diagonal n×n matrix
We say A is orthogonally diagonalizable (it is diagonalizable by an orthonormal matrix S)
Spectral Theorem: A is orthogonally diagonalziable over Rn↔AT=A
Symmetric matrices are orthogonally diagonalizable and vice versa
A=AT∧λ1=λ2→v1⋅v2=0
If A is symmetric and v1 and v2 in A are eigenvectors with different eigenvalues, then the eigenvectors must be orthogonal
In other words, any eigen vectors of different eigenspaces must be orthogonal to each other. For a matrix to be orthogonally diagonalizable means that its eigenspaces are orthogonal to each other, which would allow us to find orthogonal eigenvectors A that is symmetric has n real eigenvalues if they are counted with their algebraic multiplicities
Ordinary
Abbreviated as ODE sometimes
Only depends on one variable
Uses d to denote a differential
E.g.: dt2d2x+ω2x=0
The oscillating mass equation has a search function x(t) that depends only on time t
Partial
Depends on multiple variables
Uses ∂ to denote a differential
E.g.: ∂x2∂2E+∂y2∂2E+∂z2∂2E=c21∂t2∂2E
The wave equation has a search function E(t,x,y,z) that depends on multiple variables
1st Order
Highest order derivative is one
E.g.: y′+y2=9
2nd Order
Highest order derivative is two
E.g.: y′′+y′=3t
Higher Order
Highest order derivative is greater than 2
E.g. y(4)+y′′+ey=cost, where the order of this ODE is 4
Linear
An ODE is linear if it can be written in the form of: pn(t)y(n)+…+p1(t)y′+p0(t)y=q(t)
If q(t)=0, the linear ODE equation homogeneous. Otherwise, it’s inhomogeneous
Homogeneous q(t)=0 for the linear ODE
E.g.: y′′+3ty=0
Inhomogeneous q(t)=0 for the linear ODE
E.g.: t2y′′+cos(t)y=logt
Separable
A differential equation is separable if it can be written in the form of: y′=f(t)g(y)
The next result is to write the equation as: g(y)dy=f(t)dt, which allows for the y and t terms to be on opposite sides of the equation, setting up for integration.
Differential Forms
A differential form is a multivariable expression of the form: P(x,y)dx+Q(x,y)dy
A differential form is exact if P(x,y)dx+Q(x,y)dy=dF(x,y)
i.e.: ∃F:P=∂x∂F,Q=∂y∂F
An easy way to check for exactness is if ∂y∂P=∂x∂Q
Linearly Independent/Dependent
An equation is linearly dependent on other equations if we can write yn=Cn−1yn−1+…+C1y1
An equation is linearly independent if yn cannot be written as a linear combination of other equations
Fundamental Set of Solutions
If y1,y2,…,yn are linearly independent solutions to a linear homogeneous ODE, then they are considered the fundamental set of solutions
Phase Lag form C1cos(ωt)+C2sin(ωt)=Acos(ωt−ϕ) A=C12+C22 ϕ=angle between OC and the positive x-axis,C=(C1,C2)
Solving Methods
1st Order Linear Homogeneous ODE dtdy=f(t)g(y) ∫g(y)dy=∫f(t)dtif explicit solution existsy(t)=F(t)+C
Sometimes, there will only be implicit solutions of the form G(y)=F(t)+C
1st Order Linear Inhomogeneous ODE y′+fy=g, where y, f, and g are functions of t y=u1∫ugdt, where u=e∫f(t)dt u(t) is often called the integrating factor
For the steps in between:
(1) y′+fy=g
(2) uy′+ufy=ug
(3) (uy)′=ug
(4) uy=∫ugdt
(5) y=u1∫ugdt
From (3) to (2): u=e∫f(t)dt, so (uy)′=uy′+u′y=uy′+f(t)e∫f(t)dty=uy′+ufy
Finding Values of Derivatives of Exact Differential Forms F(x,y)=C,y′(x0)=?
(1) Find y(x0) using F(x,y)=C
(2) Calculate dF and simplify using the form Pdx+Qdy=0
(3) Qdy=−Pdx
(4) dxdy=−QP
(5) y′(x0)=−Q(x0,y0)P(x0,y0)
E.g.:
Find y′(21) for x2+y2=1 when y=y(x)>0
Geometrically, this curve is a circle about the origin with a radius of one. y′(21) would correspond to the slope of the line at the point 45° between the x and y axes, which by intuition, would mean a slope of −1. Now for the calculation by implicit differentiation of this exact differential form:
(1) F(x,y)=x2+y2=1→y(21)=21
(2) ∂x∂F=2x,∂y∂F=2y→2xdx+2ydy=0
(3) 2ydy=−2xdx
(4) dxdy=−yx
(5) y′(21)=−2121=−1
Implicit Differentiation of Exact Differential Forms
If Pdx+Qdy=0 is exact (i.e. ∂x∂F=P,∂y∂F=Q), then F(x,y)=C is a solution to Pdx+Qdy=0
(1) Check ∂y∂P=dx∂Q
(2) F=∫Pdx+g(y)
(3) Q=∂y∂F=∂y∂(∫Pdx+g(y))
(4) g(y)=∫(Q−∂y∂(∫Pdx))dy
(5) F=∫Pdx+∫Qdy−∫∂y∂(∫Pdx)dy
Implicit Differentiation of Inexact Differential Forms Pdx+Qdy=0→Find F(x,y):∂x∂F=P,∂y∂F=Q
If Pdx+Qdy not exact, then multiply by μ(x) or μ(y) (sometimes, μ may be a function of both x and y, which is quite difficult) and then solve like a normal exact differential form by implicit differentiation μ(x)=e∫h(x)dx,h(x)=Q1(∂y∂P−∂x∂Q) μ(y)=e∫g(y)dy,g(y)=P1(∂x∂Q−∂y∂P)
2nd Order Linear Homogeneous ODE
(1) Find y1 and y2 as solutions to your differential equation
(2) If W(t)=0, then y1 and y2 are a fundamental set of solutions
Variation of Parameters
Given a solution y1, we try to variate y2 by a factor of u
(1) y2=uy1
(2) a(uy1)′′+b(uy1)′+c(uy1)=0
(3) u(ay1′′+by1′+cy1)+2au′y1′+au′′y1+bu′y1=0
(4) u′′=0→u′=1→u=t
2nd Order Linear Homogeneous ODE with Constant Coefficients
For ay′′+by′+cy=0: λ=2a−b±b2−4ac b2−4ac>0→y1=eλ1t,y2=eλ2t b2−4ac=0→y1=eλt,y2=teλt b2−4ac<0→y1=ertcosωt,y2=ertsinωt,λ=r±iω
2nd Order Linear Inhomogeneous ODE with Constant Coefficients y′′+py′+qy=f(t)
(1) Solve yh for homogeneous y′′+py′+qy=0
(2) Find a particular solution yp by undetermined constant coefficients or ERF
(3) If f(t) is a sum of two terms say f1(t) and f2(t) with solutions yp1 and yp2, then yp=yp1+yp2
(4) y=yh+yp
(5) Solve IVP
Undetermined Constant Coefficients
When the inhomogeneous portion of the differential equation is f(t), we make a guess/ansatz in the form of yp as indicated in the table below
| f(t) | yp |
| | |
| cntn+cn−1tn−1+…+c1t1+c0 | antn+an−1tn−1+…+a1t+a0 |
| cos(ωt) or sin(ωt) | acos(ωt)+bsin(ωt) |
| ert | p(t)ert,p(t) is a polynomial |
| ertcos(ωt)+ertsin(ωt) | aertcos(ωt)+bertsin(ωt) |
Exponential Response Functionn ay′′+by′+cy=ert yp=p(r)ert=p′(r)tert=p′′(r)t2ert
In the above equation, go the next yp from left to right if any of the p(n)(r) are 0 ay′′+by′+cy=ertcos(ωt) yp=Re(p(r+iω)e(r+iω)t) ay′′+by′+cy=ertsin(ωt) yp=ℑ(p(r+iω)e(r+iω)t)
Homoegenous Planar System
For systems of differential equations:
x_2' = cx_1 + dx_2$$
or
$$x' = ax + by\\
y' = cx + dy$$
where $x_1,x_2,x,y$ are functions of $t$
Rewritten as $x' = Ax, ~ x = \begin{bmatrix}x_1 \\ x_2\end{bmatrix}, ~ A = \begin{bmatrix}a & b \\ c & d\end{bmatrix}$
For a $2\times2 ~ A$, the fundamental set of solutions consists of two equations
$\lambda_1,\lambda_2 \in \mathbb{R},\lambda, \bar\lambda \in \mathbb{C}, v_1,v_2,v,\bar v = \text{eigenvectors}$
$\lambda = r + i\omega$
$\bar\lambda = r - i\omega$
$v = v_1 + iv_2$
$\bar v = v_1 - iv_2$
$v$ and $\bar v$ are eigenvectors for $\lambda$ and $\bar \lambda$ respectively
Two distinct real eigenvalues:
$x(t) = C_1e^{\lambda_1 t}v_1 + C_2e^{\lambda_2 t}v_2$
- Phase Portraits



Complex eigenvalues:
$x(t) = C_1\frac{e^{\lambda t}v + e^{\bar\lambda t}\bar v}{2} + C_2 \frac{e^{\lambda t}v - e^{\bar \lambda t}\bar v}{2i} = C_1\mathbb{R}e(e^{\lambda t}v) + C_2\Im(e^{\lambda t}v)$
- Phase Portraits



Repeated real eigvenvalues and two linearly independent eigenvectors
$x(t) = C_1e^{\lambda t}v_1 + C_2e^{\lambda t}v_2$
Repeated real eigvenvalues and only one linearly independent eigenvector
$x(t) = C_1e^{\lambda t}v_1 + C_2e^{\lambda t}(tv_1 + v_2)$
$v_1 = \text{eigenvector for }\lambda$
$v_2 = \text{generalized eigenvector for } \lambda : (A - \lambda I_2) \cdot v = v_1,\text{ where }v_2 \text{ is the non }t\text{-component of }v$
Phase Portraits
Sketches of solutions to x′=Ax(t), or the planar system as time progresses
As t→∞, a planar system is
Stable: If all solutions tend to 0
Unstable: If one or more solutions are unbounded
Semi-stable: If all solutions are bounded but do not tend to 0
Matrix Exponential
e^{\lambda_1 t} & \cdots & 0\\
\vdots & \ddots & \vdots\\
0 & \cdots & e^{\lambda_n t}
\end{bmatrix}S^{-1}$$
If $A$ not diagonalizable:
$e^{At} = \Phi(t)\Phi^{-1}(0)$
$$\Phi = \begin{bmatrix}
\vec x_1 & \vec x_2
\end{bmatrix}$$
$\vec x_1$ and $\vec x_2$ are solutions to repeated real eigenvalues with only one independent eigenvector case of planar system for $A$
Inhomogeneous Planar System
(1) x′=Ax(t)+r(t)
(2) x=xh+xp
(3) xh=x for x′=Ax
(5) u=∫Φ−1rdt,Φ=[x1x2]
(4) xp=Φ⋅u A diagonalizable: x(t)=Sy(t) y′=Dy+S−1r(t)
Theorems
∃R:(t0,y0)∈R:f(t0,y0) is defined and continuous→∃y(t)
Existence theorem for an equation y′=f(t,y) with an IVP
∃R:∂y∂f is defined and continuous on R, a rectangle:(t0,y0)∈R→∃!y(t)
Uniqueness theorem for y′=f(t,y) with IVP
Superposition Principle
If y1,y2,…yn are solutions to a linear homogeneous differential equation, then so are i=1∑nCiyi for Ci∈R (i.e. their linear combination is also a solution)
Operations in Polar Form
(1) Addition/Subtraction: Convert to a+bi form
(2) Multiplication: (r1eiθ1)(r2eiθ2=r1r2ei(θ1θ2)
(3) Division: r2eiθ2r1eiθ1=r2r1ei(θ1−θ2)
(4) Conjugation: zˉ=re−iθ
Rules of Operations z,z1,z2,w∈C:
(1) Commutativity of addition: z+w=w+z
(2) Commutativity of multiplication: zw=wz
(3) Associativity of addition: (z1+z2)+z3=z1+(z2+z3)
(4) Associativity of multiplication: (z1z2)z3=z1(z2z3)
(5) Distributivity: w(z1+z2)=wz1+wz2 sin(θ)=2ieiθ−e−iθ cos(θ)=2eiθ+e−iθ
Ceikθ=C(cos(kθ)+isin(kθ))
This is called Euler’s Formula
Explanation
The intuition is that ex=dxdex, and so if an object’s position with respect to time were modeled by x(t)=et, its velocity would be represented by v(t)=et.
Now when multiplying a number by i on the complex plane, it has the effect of rotating the number by 90 degrees counter clockwise.
E.g. i⋅(a+bi)=−b+ai, where we went from (x,y)→(−y,x)
Graph showing this: https://www.desmos.com/calculator/cg4ipkcwz3
So then dtdeit=ieit, which means the velocity vector our object always points 90 degrees counter clockwise from the position vector.
Starting at t=0, x(0)=ei⋅0=1, meaning the velocity vector is i⋅1=i, which is indeed 90 degrees counter clockwise from the position vector.
From these initial conditions and the fact that the velocity vector always points 90 degrees counter clockwise from the position vector, it should be apparent now that the path of our object should trace out a circle.
Assuming constant speed of 1 (∣v∣=1), then it takes 2π units of time for the circle to be fully traced out. Notice now that if we treat the complex plane as a Cartersian plane, the parameterization of this circlular path is (cos(t),sin(t)), which directly translates to Euler’s Formula above.
Video explaining this: https://www.youtube.com/watch?v=v0YEaeIClKY eiπ=−1 is Euler’s Identity
Ln(z)=ln∣z∣+iθ,−π<θ≤π
Called the principal or complex log
(1) eLn z=z
(2) Ln z=lnz when z is a positive real number
(3) Ln(z1z2)=Ln(z1)+Ln(z2)+2πin,n={−1,0,1}
(4) Ln(z2z1)=Log(z1)−Log(z2)+2πin,n={−1,0,1} θ1+θ2>π→n=−1 θ1+θ2≤−π→n=1 −π<θ1+θ2≤π→n=0
Fourier Transform 〜
g^(f)=∫t1t2g(t)e−2πiftdt https://www.youtube.com/watch?v=spUNpyF58BY
This function takes a frequency f, and from g(t), your signal function (e.g. your signal composed of many different pure signals), determines the likelihood that that frequency is an original pure component of the total signal
As time t2−t1 increases, g^(f) becomes more defined, as it means that there is more time to distill the signals/clear uncertainty
The e−2πif part can be thought of as a wrapping g(t) around a circle in the imaginary plane in the clockwise direction (due to the negative sign) g^(f) is akin to the center of mass of a curve as g(t) is wrapped around your circle
Uncertainty ± Error propagation in addition or subtraction Q=a+b+…+c−(x+y+z) δQ=(δa)2+(δb)2+…+(δc)2+(δx)2+(δy)2+…+(δz)2+ Error propagation in multiplication or division Q=xy…zab…c δQ=∣Q∣(aδa)2+(bδb)2+…+(cδc)2+(xδx)2+(yδy)2+…+(zδz)2 Error propogation in multiplication or division by a constant Q=Cx δQ=∣C∣∣x∣δx Error propogation in exponentiation Q=xn δQ=∣Q∣n∣x∣δx
Constants ε0 g=9.81s2m (Acceleration due to gravity on Earth) G=6.6743⋅10−11kg⋅s2m3 (Universal gravitational constant) e−=−1.60217663⋅10−19C (Charge of electron) k=4πϵ01=8.9875517923⋅109C2Nm2 (Coloumb’s constant) ε0=8.85418782⋅10−12Nm2C2 (Permitivity of freespace, electric constant)
F=9.6485538⋅104mole−C (Faraday’s constant)
This is the charge in Coloumbs of a mole of electrons μ0=4π⋅10−7AT⋅m=1.25663706⋅10−6AT⋅m (Magnetic Constant)
Newtonian Mechanics 🚙 F=ma Fg=mg,g=9.8sm on surface of Earth FG=Gr2m1m2,G=6.6743⋅10−11kg⋅s2m3 Fnormal force=N, points perpendicularly from surface
Fkinetic friction=fk=μkN μk is the coeffecient of kinetic friction of the surface fk≤fs p=mv F=ΔtΔp J=Δp=FΔt=∫t0tfF(t)dt (for non-constant force) p0=pf KE0=KEf for elastic collisions
Rotational Mechanics 🎡 ac=rv2 I=∫r2dm I=Icenter of mass+MD2 I=Iring=MR2 Idisk=21MR2 Irod or plank=121ML2 Isphere=52MR2 x=rθ,v=rω,a=rα α=dtdω=dt2d2θ
τ=r×F
i.e. ∣τ∣=∣r∣∣F∣sinθ, where θ is the angle between the radius vector (pointing away from rotation axix) and the force vector
Torque follows the right-hand rule, where the forefinger first points in the direction of the radius vector and then curls to point in the direction of the force vector, resulting in an extended thumb pointing in the direction of the torque
Energy ⚡️ W=∫x0xfF⋅dx ΣWby conservative forces=ΔK ΣWby non-conservative forces=ΔE Wby conservative forces=−ΔU ΔE=ΔK+ΔU(=0 when there are no non-conservative forces) E=K+U K=21mv2
U=−∫refrF⋅dr,F=−dxdU ref can be defined to be some arbitrary reference point relative to x. If ref is consistent between U0 and Uf, ΔU will still be the same no matter the choice of ref
Pin=Pout
Pascal’s Principle: pressure exerted by fluid on a container is constant AinFin=AoutFout
P=P0+ρgΔh P is the pressure in the liquid at a certain height/of the manometer/barometer, P0 is the initial pressure (often atmospheric pressure), ρ is the density of the fluid, and Δh is the difference in height between the two menisci or between two points of interest 760torr=1atm
Manometer:
Barometer:
Fbuoyant=mg=ρVg m is the mass of the dispalced fluid, ρ is the density of the fluid displaced, V is the volume of the fluid displaced
Fr=−mkvn Fr is the retarding force on a moving object through air and is proportional to some power of the velocity of the object. This law is just an approximation though m is the mass of the object k is a constant dependent on mass that cancels out the effects of m2 \qquad v \le \text{770 mph}
\end{cases}$$
Periodic Motion 🌊
x¨+ω2x=c
SHM equation. The second time derivative of an object’s position is equal to the negative of its position times the square of the angular frequency ωspring=mk ωpendulum=Lg
ωphysical pendulum=ImgL
Simplifies to some form of Lg when I is just MR2 L is the distance from the pivot point to the center of mass of the entire physical system
ωtorsional pendulum=Iκ κ is the torsion constant and has units of radN⋅m
T=ω2π T is the time period of an oscillation f=2πω x(t)=C1eikx+C2e−ikx=D1cos(ωt)+D2sin(ωt)=Acos(ωt+ψ) v(t)=−Aωsin(ωt+ψ)=−vmaxsin(ωt+ψ) a(t)=−Aω2cos(ωt+ψ)=−amaxcos(ωt+ψ) vmax=Aω amax=Aω2
x¨+2βx˙+ω02x=x¨+mbx˙+mk=0
Damped Harmonic Oscillation differential equation ω0 is the natural frequency, or ω from the simple harmonic oscillator β=2mb, where b is the resistive constant/damping constant and m is the mass ω=w02−β2,ω0=mk,β=2mb
Underdamping: β<ω0orb<4mk (oscillation) x(t)=C1e−βtcos(ωdt)+C2e−βtsin(ωdt)=Acos(ωdt−ϕ)
Overdamping: β>ω0orb>4mk (no oscillation) x(t)=C+eS+t+C−eS−t=Ae−βtcos(wt+ϕ),S±=−β±β2−ω02
Critical damping: β=ω0orb=4mk (no oscillation) x(t)=C1eβt+C2te−βt A=x02+(ωv0+βx0)2=C12+C22 ϕ=tan−1(−ωx0v0+βx0)=angle between OC and the x-axis,C=(C1,C2)
x¨+2βx˙+ω02x=mF0cosωt
Driven oscillation differential equation β is the damping/resistivity constant ω0 is the angular frequency F is the force supplied by the driven oscillator x=xh+xp xp=Dcos(ωt−δ),D=m(ω02−ω2)2+4β2ω2F0,δ=tan−1(ω02−ω22βω)
ωR=ω02−2β2
The resonance frequency of a driven oscillation is given by the oscillator’s natural frequency and its damping coefficient ωR is specifically for the amplitude resonance
ωE=ω0 ωE is the kinetic energy resonance frequency, the frequency of the oscillator that minimizes energy loss due to the
Q=2βωR≅Δωω0 Q is the quality or Q factor of an oscillation Δω is the width of the frequency such that the square of the amplitude has decreased by half. The approximation is valid for β≪ω0
Tells how underdamped an oscillator is
As Q increases, the oscillator becomes more underdamped/less damped
Resonance: ω0>2β
No resonance: ω0<2β
Overdamping: ω0<β
0, \quad t < t_0\\
a, \quad t > t_0
\end{cases}$$
$$I(t_0, t_1) = H(t_0) - H(t_1) = \begin{cases}
0, \quad t < t_0\\
a, \quad t_0 < t < t_1\\
0, \quad t > t_1
\end{cases}$$
For $x(t') = 0, \dot x(t') = 0$
$$G(t, t') = \begin{cases}
\frac{1}{m\omega_1}e^{-\beta(t - t')}\sin (\omega_1(t - t')) \quad t \ge t'\\
0, \quad t < t'
\end{cases}$$
$x_p(t) = \int_{-\infty}^t F(t')G(t, t') dt'$
xR(x,t)=2Acos(2ψ)sin(kx−ωt+2ψ) xR is the resultant wave’s position of two interfering waves AR=2Acos(2ϕ)
Derived from the trig identities: cosα+cosβ=2cos(2α+β)cos(2α−β) sinα+sinβ=2sin(2α+β)cos(2α−β)
fobs=fsource(vsound+vsourcevsound+vobs)
Doppler effect
Observer to source is positive, vsound=343sm
θsonic boom=sin−1(vobjectvsound) θsonic boom is the angle between the line parallel to the direction of travel and the resulting shock cone vmechanical≈resistive inertiarestoring force
vstring=μT=mTL T is tension, μ is linear mass density
v_{\text{sound in air}}=\sqrt{\frac{B}{\rho}}=331K^{\degree}\sqrt{\frac{T}{273K^{\degree}}}$$v_{\text{sound}}=\sqrt{\frac{B}{\rho}}=331K^{\degree}\sqrt{\frac{T}{273K^{\degree}}} B is the bulk modulus, ρ is the volume mass density, and T is the air temperature E∝A2
P=21μA2ω2v
This is the power, P, of an oscillating string
I=AP=2ρvpmax2 I is the intensity, P is power, pmax is the pressure amplitude, ρ is the volume mass density, and v is the velocity of the wave
β=10log(I0I),I0=10−12m2W β is the decibel value, and I0 is the threshold of human hearing. W is watts fbeat=∣fa−fb∣ λlike=n2L,n∈N λmixed=n4L,n∈Nodd
∇2F=v21∂t2∂2F
This is the differential equation that describes a wave in 3D
F=F(ϕ),ϕ=ωt−k⋅r vphase=kω
Electric Charge and Electric Fields 🔋 Fe=r2kq1q1=qE,k=4πε01=9∗109C2Nm2,ε0=8.854∗10−12Nm2C2 Fe=ker2qQr=qE Fq=keqi=1∑N(ri2qi)r^i
E=kei=1∑N(ri2qi)r^i=ke∫r2dqr^
Units are CN, or mV λ=dldq,σ=dAdq,ρ=dVdq
Ering=k(z2+R2)3/22πλRz=(z2+R2)3/2Qz R is the radius of the ring z is the distance from the center of the ring to the point
Edisk=k(2πσz)(z1−z2+R21) R is the radius of the ring z is the distance from the center of the ring to the point
Φe=∮E⋅dA=ε0qen qen is the charge enclosed by the Gaussian surface dA points outward from the enclosed surface
Ω=r2A,a sphere has 4π sterdians
Akin to angles, but in 3D: θ=rl, a circle has 2π radians
\begin{cases}
0,\quad r<R\\
k\frac{Q}{r^2},\quad r\ge R
\end{cases}$$
$$E_\text{uniformly-charged sphere}=\begin{cases}
,\quad r<R
\\k\frac{Q}{r^2},\quad r\ge R
\end{cases}$$
Electric Potential ⚡️ UE=kerqQ=−∫F⋅dr=qV=21kr2Q=2ε0E2Volume,F=−∇U
VE=kerQ=kei=1∑Nriqi=k∫rdQ=−∫E⋅dr=qUE,E=−∇V
Units are V(volts)=CJ W=21∫V(r)dq VolumeUE=2ε0E2 Vconducting sphere or uniformly-charged spherical shell=krQ
p=qd p is the dipole moment, and d is the distance vector point from negative to positive charge V=kr2p⋅r^
Emax air=3⋅106CN
This is the maximum electric field that can be passed through air before sparks are formed. This process is called electric breakdown, and it can happen to any insulator. The insulator effectively becomes a conductor when such a great electric field forces it to change.
Capacitance 📸
Q=CV,C=VQ,V=CQ
Units of capacitance C are Farad (VC)
The capacitance is dependent only on the geometry of the capacitor, and not the voltage or charge held. UC=21CV2=21QV=2CQ2
Cplates=κϵ0dA κ is the dielectric constant ϵ0=4πk1=8.85418782⋅10−12N⋅m2C2(mF) A is the area of a single plate d is the distance between the plates
I=dtdQ≈ΔtΔQ=nevdA
Units are Ampere (A=sC)
Current measures the number of electrons that pass through a point per unit time
Current is not a vector because for a given bending wire, it may travel in multiple directions. n is the number of charged particles per unit volume e is the magnitude of the charge on each charged particle vd is the drift velocity of electrons, ≈10−4sm in a wire A is the area of the wire
Looks like I= Nevada
J=dAdI=n∣q∣vd J is volume current density
We don’t call it area current density because current is already a one-dimensional (rather than zero-dimensional, like charge), and so when multiplying J by the 2D quantity of area, the value essentially becomes 1D current.
E=ρJ
This is Ohm’s Law ρ is the resistivity of a material ρ=JE
σ=ρ1 σ is the conductivity of the material
ρ(T)=ρ0[1+α(T−T0)] ρ0 is the reference resistivity at a reference temperature T0, and α is the temperature coefficient of resistivity
R=AρL
Units are Ohm’s (Ω)
Resitivity is kind of like the density of people in a crouded hallway while resistance is like the total number of collisions a student may make trying to get to the end of the hall R(T)=R0[1+α(T−T0)]
V=IR
Is not necessarily true for all systems
Is more like an approximation used by engineers P=IV=I2R=RV2 I is constant in series circuits Rseries=i=0∑NRi V is constant in parallel circuits Rparallel=(i=0∑NRi1)−1
∑Iin=∑Iout
This is Kirchhoff’s Current/Junction Rule
(∑ΔV)loop=0
This is Kirchoff’s Loop Rule
Let l be the path vector of the loop Vresistor=−(I⋅l)R Vbattery=ε τRC=RC,Qmax=Cε,Imax=Rε,V(t)=CQ(t) Qcharge(t)=Qmax(1−e−τt) Icharge(t)=Imaxe−τt Qdischarge(t)=Qmaxe−τt Idischarge(t)=−Imaxe−τt
Magnetic Fields and Forces 🧲
F=q(E+v×B) B is symbol for the magnetic field and has units of Teslas (T), which is equal to A⋅mN
Gauss (G)=10−4T
Earth’s magnetic field is about 25-65 μT, or 0.25-0.65 G
Also known as Lorentz Force, where qE is the electric component of the force and v×B is the magnetic component of the force
For the force on current-carrying wire, an equivalent formula is: FB=IL×B
The magnitude of the magnetic force is then is ∣FB∣=∣q∣vBsinϕ
∮B⋅dl=μ0Ienclosed
This is known as Ampere’s Law, and it says that the vector line integral of the magnetic field around an enclosed area is proportional to the enclosed current
i.e. the summated curls of the magnetic field around a point is proportional to the current flowing through that point
Closed path integrals are evaluated counter clockwise when current points out of the page
ΦB=∬B⋅dA
Magnetic flux has units of Webers (Wb) and equals Tm2
∮B⋅dA=0
Says that no magnetic charge or magnetic monopoles exist
r=qBmv
This is the radius of a charged particle in a cyclotron/mass spectrometer and has units of meters when kg, sm, C, and T
are used for m, v, q, B respectively.
ω=mqB,T=∣q∣B2πm,f=2πm∣q∣B
The angular velocity, time period, and frequency of a charged particle circuiting in a magnetic field
v=BE
The velocity that is screened for in a velocity selector is dependent only on E and B
FB=nevdALB=IL×B
Known as the Lorentz Force L points in the direction of the current nevdA is the same as current FB=∫Idl×B
μ=IA
Known as the magnetic moment, and anything that produces one can be called a magnetic dipole
For a motor, A points perpendicular to the plane of the motor coils in the direction of thumb when curling the right hand along the direction of current flow
τ=NIABsinϕ=Nμ×B
The torque of a motor depends on the number of motor coils N, the current I, the area of the motor coils A, the magnetic field B, and the angle between the area vector and the magnetic field vector ϕ
The torque has the effect of aligning the area vector/magnetic moment with the magnetic field vector
U=−μ⋅B
The potential energy of a magnetic moment in a magnetic field
Magnets
Magnets generate magnetic fields out of their north end and into their south end
A magnetic dipole near a nonuniform magnetic field-generating magnet (e.g. a bar magnet) will be attracted if its magnetic moment points parallel with the generated magnetic field. If they point in opposite directions, they will repel each other
Magnetic dipoles that are attracted will also experience a stretching force, whereas magnetic dipoles that are repelled will experience a compression force
Hall Effect
Experiment showed that charge carriers are negative nq=−EzJxBy
B=4πμ0∣r∣2qv×r^=4πμ0∣r∣3qv×r μ0=4π⋅10−7AT⋅m=1.25663706⋅10−6AT⋅m (Magnetic Constant) r points from the moving charge to the point of B measurement r^ is the normalized (magnitued of 1) unit vector of r
B=4πμ0∫∣r∣2Idl×r^=4πμ0∫∣r∣3Idl×r
Called the Biot-Savart Law r and r^ are the same as in the above equation
∣dB∣=4πr2μ0Idlsinϕ
Works when ∣dB∣ and dl are sufficiently small ϕ is the angle between the magnetic field vector and the radial unit vector, which points from the current component to the point of measurement
Bwire=2πrμ0I
The direction of the B can be found by the right thumb to point in the direction of I and then curling the other four fingers, which will point in the direction of the generated B
Bwire length=4πμ0Irr2+(2ℓ)2ℓ
This is the magnetic field at a point a distance r from the conducting wire of length ℓ
Bcoil=2(h2+R2)3/2μ0NIR2
The direction of B points in the direction of the right thumb when the other four fingers curl in the direction of I in the coil N is the number of loops R is the radius of the coil h is the distance between the center of the coil and the point of measurement
Bsolenoid=ℓμ0NI
This is the magnitude of the magnetic field inside a solenoid N is the number of coil turns ℓ is the length of the solenoid
The solenoid points in the direction of the right thumb when curling the fingers in the direction of current
Btoroidal solenoid=2πrμ0NI
This is the magnitude of the magnetic field inside a toroidal solenoid N is the number of wire loopings around the toroid
Binside wire=2πR2μ0Ir r is the distance from the center of the wire to the radius of measurement R is the radius of the wire
Bparallel wires2πrμ0I1I2
Wires with current flowing in the same direction will attract each other while wires with current flowing in opposite directions will repel each other
Bfrom μ=2π(x2+a2)3/2μ0μtotal x is the distance from the magnet to the point of interest a is the width of the magnet μtotal is often found from M, the magnetization of a magnet
Quantum Magnetism ⚛️
μ=ml2meℏ μ is the magnetic moment of an atom with an electron orbiting it (Bohr magneton) ml is the magnetic quantum number e is the charge of the electron ℏ=2πh=1.05457266⋅10−34J⋅s (Reduced Planck’s constant) m is the mass μB=2meℏ=9.274⋅10−24A⋅m2
M=Vμtotal M is the magnetization of a material and is the net magnetic moment per unit volume μtotal is often 0 in atoms because the orbital and spin magnetic moments of electrons tend to cancel each other out in atoms
For atoms like iron where there are free unparied d orbital electrons, they can achieve a net magnetic moment
Unpaired e− = paramagnetic
Paired e− = diamagnetic
B=B0+μ0M
The overall magnetic field inside a material is equal to the external magnetic field applied plus the permeability of free space times the magnetization of the material
If M is aligned/parallel to B0, then the material is paramagnetic. Otherwise, it’s diamagnetic
So B is greater in paramagnetic materials, meaning they are more magnetic (e.g. iron is more magnetic and it’s paramagnetic)
Km=B0B=μ0μ
Called the relative permeability of a material Km=1 means the material is like a vacuum in terms of the strength of the magnetic field inside it Km>1 means the material is paramagnetic, or has a stronger magnetic field inside than in a vacuum for the same external B0 Km<1 means the material is diamagnetic, or has a weaker magnetic field inside than in a vacuum for the same external B0
χm=Km−1
Called the magnetic susceptibility of a material
Common magnetic susceptibilities:
Hysteresis
Superconductors
Induction 🛵
∣ε∣=N∣dtdΦB∣
Called Faraday’s Law of Induction ε is the emf in a closed loop N is the number of loops of coil dtdΦB is the change in magnetic flux through the area bound by the closed loop
Units are sWb, and Wb=Tm2
Lenz’s Law
States the direction of ε in Faraday’s law
The direction is that of the four right fingers curled when the thumb points in the direction of decreasing magnetic flux ε=∮(v×B)⋅dl
∣εalternator∣=ωBAsinωt ω is the angular velocity of the alternator rotation
∣εalternator average∣=π2NωBA N is the number of loops of the alternator coil ω is the angular velocity of the alternator
∣εslidewire∣=vLB v is the velocity of the rod L is the length of the rod B is the magnitude of the magnetic field
∣εfaraday disk dynamo∣=21ωBR2 ω is the angular velocity of the dynamo R is the radius of the disk
M=I1N2ΦB2=I2N1ΦB1 M is the mutual inductance between two solenoids on the same axis M is also the constant of proportionality between ΦB2 and I1 or between ΦB1 and I2
Units are Henry (H)=AWb
L=INΦB L is the self-inductance, which is the constant of proportionality between ΦB and I through the solenoid
Units are also H
ε=−MdtdI=−LdtdI
Inductors resist the flow of current. i.e. when you first switch on a circuit, the inductor induces an emf in the opposite direction to the circuit flowing through it. Over time, the current flow will steady out, meaning dtdI=0. This steady current means the inductor won’t induce any emf. When you switch the circuit off, the current decreases, meaning the inductor will induce an emf in the direction that circuit is flowing.
The water-pipe analogy for inductors are water wheels. It takes time/energy to get a water wheel going, but once it’s going, it keeps going even when there is nothing driving the water forward. Uinductor=21LI2
u=2μ0B2 u is the magnetic energy density Lseries=i=0∑NLi Lparallel=(i=0∑NLi1)−1 τRL=RL,Imax=Rε,Vdrop=−ε=LdtdI Icharge=Imax(1−e−τt) Idischarge=Imaxe−τt ωLC=LC1,Qmax=Cε,Vcapacitor(t)=CQ(t),Vinductor(t)=−εinductor=LdtdI Q(t)=Qmaxcos(ωLCt+ϕ) I(t)=−ωQmaxsin(ωLCt+ϕ) U=2CQmax2=UB+UE=21LI2+2CQ2 β=2LR,ωLRC=ωLC2−β2 Q(t)=Qmaxe−βtcos(ωLRCt+ϕ)
Under damping: β<ω0orR2<C4L (oscillation)
Critical damping: β=ω0orR2=C4L (no oscillation)
Over damping: β>ω0orR2>C4L (no oscillation)
AC Current ∿ v(t)=Vcos(ωt),i(t)=Icos(ωt) Vrms=2V,Irms=2I vR(t)=IRcos(ωt) vL(t)=IXLcos(ωt+2π),XL=ωL vC(t)=IXCcos(ωt−2π),XC=ωC1
vLRC=IZcos(ωt+ϕ),Z=R2+(XL−XC)2,ϕ=tan−1(RXL−XC) Z is the impedence of an LRC circuit, which is a sort of resistance of the circuit ϕ is the power factor of your LRC circuit, which is a factor of how much power is dissipated by the circuit. pR(t)=I2Rcos2(ωt) pL(t)=−I2ωLsin(ωt)cos(ωt) pC(t)=ωCI2sin(ωt)cos(ωt)
Paverage LRC=21VIcos(ϕ)
When ϕ=0 with power factor cos(ϕ)=1, you have the power delivered as P=VrmsIrms, which is the same as the power delieverd from a resistor
When ϕ=±90° with power factor cos(ϕ)=0, you have the power delivered as P=0, which is the same as the power delieverd from a pure inductor or capacitor
ω0=LC1
The resonance frequency of an AC LRC circuit is the same as the resonance frequency of a DC LC circuit
V1V2=N1N2
For a transformer, the voltage received by a secondary winding is proportional to the ratio of the windings (so more windings on the second winding means greater output voltage) and to the voltage through the first winding N2>N1→V2>V1, which is a step-up transformer N2>N1→V2>V1, which is a step-down transformer R=(N1N2)2I1V1
Maxwell’s Equations 📡
∮SE⋅dA=ε0qenc
Gauss’s Law
States that the electric flux through an enclosed surface is proportional to the charge that the surface encloses
∮SB⋅dA=0
Gauss’s Law for magnetism
States that the magnetic flux through an enclosed surface is equal to 0
i.e. there is no such thing as a magnetic dipole/charge
∮CE⋅dl=−dtdΦB
Faraday’s Law
States that the circulation/curl of an electric field about a loop scales with changing magnetic flux through the loop
Direction of current is dictated by Lenz’z law
∮CB⋅dl=μ0(Ic+εdtdΦE)enc
Ampere’s Law
States that the curl of a magnetic field around a loop scales with the current and displacement current (changing electric field) through the loop
Displacement current density is εdtdE
∇⋅E=ε0ρ
The divergence of the electric field at a point is proportional the point’s charge density
∇⋅B=0
The divergence of the magnetic field at a point is zero; there are no magnetic sources or sinks (monopoles) ∇×E=−∂t∂B ∇×B=μ0(J+ε0∂t∂E)
Electromagnetic Waves 🌈 c=λν=μ0ε01
E×B=cprop cprop is the direction of propagation of an EM wave
\Vert\vec B\Vert$$
- $E(x,t)=E_\text{max}\cos(kx-\omega t)$
If $\vec c$ goes in the same direction as positive $x$, then $\omega$ is positive. Otherwise, $\omega$ is negative
- $B(x,t)=B_\text{max}\cos(kx-\omega t)$
Same rules apply as to $E(x,t)$
$k=\frac{2\pi}{\lambda},\space \omega=\frac{2\pi}{T}=2\pi f$
$\lambda_{\text{like}}=\frac{2L}{n},\space n\in \mathbb{N}$
$\lambda_{\text{mixed}}=\frac{4L}{n},\space n\in \mathbb{N_{odd}}$
- $n=\frac{c}{v}=\sqrt{\kappa\kappa_\mu}$
The index of refraction
Ratio of the speed of light in a vacuum to the speed of light in a medium
$\kappa$ is the relative permitivity or dielectric constant of the medium
$\kappa_\mu$ is the relative magnetic permeability of the medium
$u=\frac{I}{c}=\frac{1}{2}\varepsilon_0E^2+\frac{1}{2\mu_0}B^2=\varepsilon_0E^2=\frac{B^2}{\mu_0}$
- $\vec S=\frac{1}{\mu_0}\vec E\times\vec B=\frac{E_\text{max}B_\text{max}}{\mu_0}\cos^2(kx-\omega t)$
Poynting vector
Average of $S$ is light intensity, or the power per unit area
$P=\oint\vec S\cdot d\vec A$
- $I=S_\text{average}=\frac{E_\text{max}B_\text{max}}{2\mu_0}=\frac{1}{2}\varepsilon_0cE^2_\text{max}$
Has units of $\frac{W}{m^2}$, or power per unit area
- $p_\text{rad absorbed}=\frac{I}{c}$
Defined as radiation pressure
Same units as regular pressure
$p_\text{rad reflected}=\frac{2I}{c}$
$E_\text{standing}(x,t)=-2E_\text{max}\sin(kx)\sin(\omega t)$
$B_\text{standing}(x,t)=-2B_\text{max}\cos(kx)\cos(\omega t)$
$x_{\text{nodal planes of }\vec E}=n\frac{\lambda}{2};\space n\in \mathbb{N}$
$x_{\text{nodal planes of }\vec B}=n\frac{\lambda}{4};\space n\in \mathbb{N}_\text{odd}$
θr=θa
The angle between the reflected light ray and the normal of the surface is equal to the angle between the incident light ray and the normal of the surface
nasin(θa)=nbsin(θb)
Snell’s Law θ is measured as between the light ray and the normal to the medium surfaces na is the index of refraction for the medium the light enters from nb is the index of refraction for the medium the light enters
Other forms are: sin(θb)sin(θa)=vbva=nanb
λ=nλ0,k=nk0 λ is the wavelength of light in the medium λ0 is the wavelength of light in a vacuum n is the index of refraction k is the wave number of light in the medium k0 is the wave number of light in the vacuum f stays constant between mediums, λ shortens to compensate for the wave slowing down
θcrit=sin−1(nanb)
When an light ray makes an angle θcrit with the normal of the medium surfaces, the resulting light ray travels along the seam between the mediums rather than reflecting or refracting na is the medium that the light enters from nb is the medium that the light enters
n∝λ1
For most materials, this is the case
For many materials, the amount of refraction depends on the wavelength of light. If this is true, this dependence is called dispersion
I=Imaxcos2(ϕ),Imax=2I0
Malus’s Law
The intensity through a second polarizer depends on the angle ϕ between the two polarizers’ axes
tan(θp)=nanb
Brewster’s Law
The angle of incident light at which none of the reflected light has light polarized in the parallel to the plane of incidence
A consequence of this law is that the reflected and refracted ray are perpendicular to eachother θp is the angle between the polarized light and the
Irayleigh scattered∝λ41
The intensity of scatttered light due to Rayleigh scattering is inversely proportional to the fourth power of the wavelength of light. So blue light for example scatters much more than red light
m=yy′=−ss′=−nbsnas′
The lateral magnification of an image m<0 corresponds to an inverted image y′ is the image height y is the object height s is the distance from the object to the mirror surface s′ is the distance from the object to the mirror surface s>0 when the object’s on the same side as the incident ray s′>0 when the image’s on the same side as the reflected or refracted ray
mtotal=m1m2…mn
The total magnification of a series of lens is the product of each lens’ magnifications
s1+s′1=f1=R2
Image position for a spherical mirror of radius R f is the focal length of the spherical mirror
Same sign conventions as in above equation R>0 when the center of the spherical mirror is on the same side as the reflected light
Essentially assumes R≪s
sna+s′nb=Rnb−na
Spherical lens equation for refraction between two media with different refractive indexes na is the refractive index of the medium that the incoming light first passes through nb is the refractive index of the medium that the refracted light passes through
Same sign rules as above equation
Essentially assumes R≪s
f1=(n−1)(R11−R21)
Thin lensmaker equation f is the focal length
Assumes the lens’ thickness is less than the curvature of the lens (d≪R)
f1=(n−1)(R11−R21+nR1R2(n−1)d)
Full lensmaker equation
| s>0 | s′>0 | R>0 | |
| ------------------------------ | ---------------- | --------------- | ----------------------------- |
| If the X is on the same side as Y light rays | object; incoming | image; outgoing | center of curvature; outgoing |
f-number=Df
The f-number of a camera depends on the focal length f and the diameter of the aperature D f-numbers are generally given in the format of f-numberf. For example, 1.4f would correspond to an f-number of 1.4 I∝(f-number)21
P=f1 P is the optical power of the lens f is the focal length of the lens and =2R, where R is the radius of the lens
The optical power of a lens determines how much bends light toward a point P>0 means f>0 and that the lenses are converging for hyperopia (far-sighted, short eyeballs) P<0 means f<0 and that the lenses are diverging for myopia (near-sighted, long eyeballs)
Increasing ∣P∣ means lowering f and R, which means more converging or diverging of light
| | Nearsighted | Farsighted |
| --------------------------- | ------------------------------------------- | ------------------------------------------- |
| Can see | Near things | Far things |
| Can’t see | Far things | Near things |
| X-opia | Myopia | Hyperopia |
| X-eyeballs | Long eyeballs | Short eyeballs |
| Lenses needed | Divergent lens | Convergent lens |
| f needed for focal lenses | f<0 | f>0 |
| P needed for focal lenses | P<0 | P>0 |
| Eye size through lens | Eyes look small from observer’s perspective | Eyes look large from observer’s perspective |
M=θθ′=fsnearpoint M is the angular magnification of θ′ is the angle subtended by the image height at the point of observation θ is the angle subtended by the object height at the point of observation snearpoint is often taken as 25 cm for the human eye f is the focal length of the lens
Assumes that the object angle is small to allow for θ′≈θ≈sinθ=fy
Mmicroscopes=m1M2=f1f2snearpoints1′
Microscopes use two converging lens: the objective lens and the eyepiece snearpoint is often taken as 25 cm for the human eye s1′ is the length between the objective lens and the focal point of the eyepiece f1 and f2 are the focal length of the lenses
Mtelescopes=−f2f1
Telescopes use two converging lens like microscopes
Interference 📡
\text{Constructive Interference: }r_2-r_1=m\lambda, \space m\in \mathbb{Z}$$\text{Destructive Interference: }r_2-r_1=(m+\frac{1}{2})\lambda, \space m\in Z ri is the distance from the point source to the point of interest λ is the wavelength of light
Constructive interference for two points produces linear nodal curves
Destructive interference for two points produces hyperbolic antinodal curves
The central bright peak corresponds to m=0 for constructive inteference
The nth bright peak corresponds to m=±n for constructive inteference
The nth dark band corresponds to m+21=±2n
r2−r1=dsinθ
The difference in distances from two double slits to the point of interest is approximated by the above equation d is the distance between the slits θ is the angle between the normal plane and the line formed between the center of the double slit and the point of interest
Assumes d≪R, the normal distance from the slits to the screen
ym=Rtanθm≈Rsinθm
For constructive interference, ym is the distance between the center of the intensity pattern and the mth bright band
For destructive interference, ym is the distance between the center of the intensity pattern and the mth dark band
The approximation assumes d≪R and ym≪R
When θm=θ0, this corresponds to the central bright peak θm=θn corresponds to the nth bright peak
I=I0cos2(2ϕ)≈I0cos2(λRπdy),I0=2ε0cE2,ϕ=λ2πdsinθ,sinθ≈Ry
The intensity of light at each point depending on either the angle θ between the center line (the normal to the surface that intersects the center of the double slit) and the point of interest or depending on the shortest height y between the center line and the point of interest
This equation is for intereference, but it does not account for diffraction. The equation predict equal amplitudes of all of the intensity peaks, but in reality, the peaks diminish as θ deviates further from 0
The approximations are valid only when d≪R NOTE: This formula and many others work only in radians mode (as indicated by the π’s)
na>nb
na=nb
na<nb
ϕreflected=0
No reflection
ϕreflected=π
- $\text{Constructive: }2t=(m+\frac{\Delta\phi}{2\pi})\lambda,\space \text{Destructive: }2t=(m+\frac{1}{2}-\frac{\Delta\phi}{2\pi})\lambda \space \\m\in0,\mathbb{Z}^+$
$t$ is the thickness of a thin film
$\lambda$ is the wavelenght of light in the medium of the thin film (may have to use $\lambda=\frac{\lambda_0}{n}$)
When only one surface undergoes a half-cycle phase shift, then $\frac{\Delta\phi}{2\pi}=\frac{1}{2}$. Otherwise, $\frac{\Delta\phi}{2\pi}=0$


$y_{m_\text{michelson inferometer}}=m\frac{\lambda}{2}$
- Diffration 🔦
- $\sin(\theta_\text{dark})=\frac{m\lambda}{a},\space m=\mathbb{Z}-\{0\}$
The angular location of the dark fringes for far-field (Fraunhofer) diffraction
The $n$th dark band would correspond to $m=\pm n$
- $y_m=R\sin\theta_m$
$\theta_m$ can be found from $\sin(\theta_\text{dark})\approx\frac{m\lambda}{a}$, which holds when the screen is far from the light obstacle
Interestingly, the central bright band is twice as large as the other bright bands
- $I=I_0\cos^2(\frac{\phi}{2})(\frac{\sin(\beta/2)}{\beta/2})^2,\space\beta=\frac{2\pi}{\lambda}a\sin(\theta)$
Intensity of light at a certain angle $\theta$ or height
The $\cos^2(\frac{\phi}{2})$ factor comes from interference of the waves
The $(\frac{\sin(\beta/2)}{\beta/2})^2$ factor comes from diffraction of light
- $I_m\approx\frac{I_0}{(m+\frac{1}{2})^2\pi^2},\space m\in\mathbb{Z}^+$
The intensity of a bright band for far-field (Fraunhofer) diffraction
The $n$th bright peak corresponds to $m=\pm n$
The intensity at $m=0$ is just $I_0$ (doesn’t appear that way from this equation)
- $\text{Missing maxima at }k\text{th peak when }d=ka$
For a double slit experiment with finite-width slits, the $k$th bright peak is missing when the slit separation distance $d$ is an integer multiple of the slit width $a$ by a factor of $k$
- $I_\text{principal peaks}=N^2I_0$
When there are $N$ slits for diffraction grating, the intensity of the principal peaks are $N^2$ times the central peak intensity for when there is one slit
$\#\text{ minima between principal peaks}=N-1$
$\text{Principal peak width}\propto\frac{1}{N}$
- $\sin\theta_1=1.22\frac{\lambda}{D}$
Light outgoing from a circular aperture produces a circular diffraction pattern with a central bright circle called the Airy disk
$\theta_1$ is the angle between the center line and the line formed from the center of the aperture and the edge of the Airy disk
$\lambda$ is the wavelength of light
$D$ is the diameter of the aperture
- Special Relativity 🚀
*NOTE*: I have only verified these equations have only been verified for inertial reference frames (no accelerating reference frames)
- Einstein’s Postulates
(1) The laws of physics are the same from all inertial reference frames
(2) The speed of light in a vacuum is a constant $c$
(3) For low speed, classical Newtonian physics should hold
- $\gamma=\frac{1}{\sqrt{1-\frac{v^2}{c^2}}}$
Called the Lorentz’s factor
Is a ubiquitous factor for equations in special relativity like for time dilation
$v$ is the magnitude of the velocity that a reference frame $S'$ is moving relative to $S$
When $v\ll c$, $\gamma\approx 1$
As $v\rightarrow c$, $\gamma$ increases significantly beyond $1$
- $\Delta t=\gamma\Delta t_0$
This equation describes the time dilation of an event that is stationary in $S'$ as observed by someone in the reference frame $S$
$\Delta t_0$ is the time measured of an event that occurs at the same point in $S'$ (called the proper time for this event)
$\Delta t$ is the time measured of an event in $S$
- $l=\frac{l_0}{\gamma}$
This equation describes the length contraction of an object that is stationary in $S'$ as observed by someone in reference frame $S$
$l_0$ is called the proper length of the object and is measured in $S'$.
$l$ is the length measured of the object in $S$
The length measured has to be parallel to the movement between $S$and $S'$
- $x'=\gamma(x-ut)$
The Lorentz-transformed $x$-coordinate $x'$ in $S'$ of $x$ in $S$
$u$ is the velocity of $S'$ relative to $S$
- $t'=\gamma(t-\frac{ux}{c^2})$
The Lorentz-transformed time $t'$ in $S'$ of $t$ in $S$
$t$ and $x$ are native to $S$
$v'_x=\frac{v_x-u}{1-\frac{uv_x}{c^2}}$
$p=\gamma mv$
$F^\parallel=\gamma^3ma$
$F^\perp=\gamma ma$
$F=\gamma m[a+\frac{\gamma^2}{c^2}v(\vec v\cdot \vec a)]$
$K=(\gamma-1)mc^2$
$E=K+mc^2=\gamma mc^2$
$E^2=(mc^2)^2+(pc)^2$
- Lorentz Transformation
Blog explaining hyperbolic trig functions: https://www.physicslog.com/blog/2019/03/lorentz-hyperbolic-rotation/
$$\begin{bmatrix}
t'\\
x'\\
y'\\
z'
\end{bmatrix} =
\begin{bmatrix}
\gamma & \gamma\beta & 0 & 0\\
\gamma\beta & \gamma & 0 & 0\\
0 & 0 & 1 & 0\\
0 & 0 & 0 & 1
\end{bmatrix}
\begin{bmatrix}
t\\
x\\
y\\
z
\end{bmatrix} =
\begin{bmatrix}
\cosh \theta & \sinh \theta & 0 & 0\\
\sinh \theta & \cosh \theta & 0 & 0\\
0 & 0 & 1 & 0\\
0 & 0 & 0 & 1
\end{bmatrix}
\begin{bmatrix}
t\\
x\\
y\\
z
\end{bmatrix}$$
$\beta = \frac{v}{c}$
$\gamma = \frac{1}{\sqrt{1 - \beta^2}} = \cosh \theta$
$\gamma\beta = \sinh \theta$
- $s^2 = (ct)^2 - \vec r^2$
$s$ is the preserved interval for Lorentz Transformations
$\vec r = (x, y, z)$
$$g_\text{Euclidean} = \begin{bmatrix}
1 & 0 & 0\\
0 & 1 & 0\\
0 & 0 & 1
\end{bmatrix}$$
$$g_\text{hyperbolic} = \begin{bmatrix}
1 & 0 & 0 & 0\\
0 & -1 & 0 & 0\\
0 & 0 & -1 & 0\\
0 & 0 & 0 & -1
\end{bmatrix}$$
$$s^2 = \begin{bmatrix}
ct & x & y & z
\end{bmatrix} g_\text{hyperbolic}
\begin{bmatrix}
ct\\
x\\
y\\
z
\end{bmatrix} = (ct)^2 - (x^2 + y^2 + z^2)$$
- $$\phi = \begin{bmatrix}
\omega / c & k_x & k_y & k_z
\end{bmatrix} g_\text{hyperbolic}
\begin{bmatrix}
ct\\
x\\
y\\
z
\end{bmatrix} = \omega t - \vec k \cdot \vec r$$
$$k = \begin{bmatrix}
\omega / c, k_x, k_y, k_z
\end{bmatrix}$$
- Mathematical Methods
https://www.physicswithelliot.com/all-notes
- Vector/Tensor Analysis
- Tensor Definition
Best contravariant/covariant explanation: https://www.youtube.com/watch?v=CliW7kSxxWU
Tensor explanation 1: https://www.youtube.com/watch?v=7c8Agf9qtfI
Tensor explanation 2: https://www.youtube.com/watch?v=ztUHlZftPlo
Tensor explanation 3: https://www.youtube.com/watch?v=nNMY02udkHw
A rank $n$ tensor in $m$ dimensional space has $n$ indicies (or coordinates), $m^n$ components, and transforms as described by a generalized version of the above transformation rule
E.g. the 3D stress tensor is a rank $2$ tensor and has $3^2$ components. Likewise, vectors in 3D are rank $1$ tensors and have $3^1$ components
Superscript represents column vectors:
$$T^i = (T^i)_{i = 1, 2, 3} =
\begin{bmatrix}
T^1\\
T^2\\
T^3
\end{bmatrix}$$
Subscript represents row vectors:
$$T_{i} = (T_i)_{i = 1, 2, 3} = \begin{bmatrix}
T_1 & T_2 & T_3
\end{bmatrix}$$
- $\varepsilon_\text{ijk}=
\begin{cases}
~~~1,\quad\text{even permutation of } \varepsilon_{ijk}\\
-1,\quad \text{odd permutation of } \varepsilon_{ijk}\\
~~~0,\quad \text{otherwise}
\end{cases}$
Called the Levi-civita symbol
$\varepsilon_{ijk}\delta_{jk}=0$
$\varepsilon_{ijk}\varepsilon_{ilm}=\delta_{jl}\delta_{km}-\delta_{jm}\delta_{kl}$
$M_{ij}N_{jk}=(MN)_{ik}$
$R_{ij}=R_{ji}^{-1}\text{ for orthogonal }R$
- $v_i'=R_{ij}v_j$
The vector components in a new coordinate system $S'$ that is rotated with respect to the original coordinate system $S$ can be represented by the orthogonal transformation of the vector
Vectors are defined to describe values that do not change despite changing coordinate systems. The components of the vector may change when the coordinate system is rotating or even moving, but the vector’s direction/magnitude do not change
- $\lVert\vec a\times \vec b\rVert=\lVert\vec a\rVert\lVert\vec b\rVert\sin\theta$
Called the vector product or cross product
The resulting vector is orthogonal to both vectors and follows the right hand rule where taking your right hand, you point toward $\vec a$, curl your hands toward $\vec b$, and the direction of your thumb when pointing out is the $\lVert\vec a\times \vec b\rVert$
(1) Anticommutative: $\vec a\times \vec b=-\vec b\times \vec a$
(2) Distributive: $\vec a\times(\vec b + \vec c)=\vec a \times \vec b+\vec a\times \vec c$
- $(\vec a\times \vec b)_i=\varepsilon_{ijk}a_jb_k$
$(\vec a\times \vec b)_i=a_2b_3-a_3b_2=\varepsilon_{ijk}a_jb_k$
- $\text{Triple Scalar Product}$
$$\vec A\vec B\vec C=\det\begin{bmatrix}
A_x & A_y & A_z\\
B_x & B_y & B_z\\
C_x & C_y & B_z
\end{bmatrix}
= \vec A \cdot (\vec B \times \vec C) = (\vec A \times \vec B) \cdot \vec C = \text{even permutations of }\vec A \vec B \vec C = -\text{(odd permutations of }\vec A \vec B \vec C)$$
- $\text{Triple Vector Product}$
$\vec A \times (\vec B \times \vec C) = (\vec A \cdot \vec C)\vec B - (\vec A \cdot \vec B)\vec C$
In words (this is how it should be memorized):
The value of a triple vector product is a linear combination of two vectors in the parenthesis; the coefficient of each vector is ht edot product of the other two; the middle vector in the triple product always had the positive sign
- $I_{ik}=\sum m(r^2\delta_{ij}-r_ir_j)=\int d^3r\rho(\vec r)(r^2\delta_{ij}-r_ir_j)$
Expanded form $=\begin{bmatrix}
\sum m(y^2+z^2)&-\sum mxy&-\sum mxz\\
-\sum mxy &\sum m(x^2+z^2)&-\sum myz\\
-\sum mxz& -\sum myz&\sum m(x^2+y^2)
\end{bmatrix}$
- $\text{Principle Axes}$
The principle axes of an object are the eigenvectors of the object’s diagonalized moment of inertia tensor matrix
- $E=\frac{1}{2}\omega_i\omega_jI_{ij}$
This mirrors the one-axis form of:
$E=\frac{1}{2}I\omega^2$
$L_i=I_{ij}\omega_j$
- $x^{i'}=x^{i'}(x^1,x^2,\ldots,x^N)$
Represents a generalized coordinate transformation
- ${v^i}'={\frac{{(\partial x^i)}'}{\partial x^k}}v^k$
Represents a generalized transformation of a vector. Note we have to sum over $k$ elements of the vector for each $i$ elements in the transformed vector where $k = i = \text{dimension of }v$
Additional fact: ${dx^i}' = \frac{\partial x^i}{\partial x^k}'dx^k$
An example would be the rotation transformation:
$$x' = x\cos\theta + y\sin\theta\\
y' = -x\sin\theta + y\cos\theta$$
Or using superscripts:
$${x^1}' = x^1\cos\theta + {x^2}\sin\theta\\
{x^2}' = -x^1\sin\theta + x^2\cos\theta$$
where ${x^1}'$ is the $x^1$-coordinate of a vector in the transformed basis
Now the vector would be represented by:
${v^1}' = \sum_{k = 1}^2 \frac{{\partial x^i}'}{\partial x^k}v^k = \frac{{\partial x^{1}}'}{\partial x^1}v^1 + \frac{{\partial x^{1}}'}{\partial x^2}v^2 = \frac{\partial}{\partial x^1}(x^1\cos\theta + {x^2}\sin\theta)v^1 + \frac{\partial}{\partial x^2}(-x^1\sin\theta + {x^2}\cos\theta)v^1= v^1\cos\theta + v^2\sin\theta$
${v^2}' = \ldots = -v^1\sin\theta + v^2\cos\theta$
- ${T^{ik}}'={\frac{{(\partial x^i)}'}{\partial x^j}}{\frac{{(\partial x^k)}'}{\partial x^l}}T^{jl}$
Generalized form for a tensor transformation
- $ds^2 = g_{ik} dx^i dx^k$
$ds$ here is an invariant distance between two points
$g_{ik}$ is the metric tensor with the property that $g_{ik}$ applied to
In Euclidean space, $g_{ik} = \delta_{ik}$, which gives that $ds^2 = \delta_{ik}dx^idx^k = dx_kdx^k$, which is the scalar product between the two points’ vectors, which indeed equals the square of the distance between the two points
- Linear Algebra
$\det(A=a_{ij})=\varepsilon_{ijk}a_{1i}a_{2j}a_{3k}$
- Differential Operators
- Rules
$\partial$ is only commutative with respect to other $\partial$ signs
E.g:
$\partial_i\partial_ju_i=\partial_j\partial_iu_i$ because $\frac{\partial^2 u_i}{\partial x_j\partial x_i}=\frac{\partial^2 u_i}{\partial x_j\partial x_i}$
However, $\partial_iu_iv_k\ne \partial_iv_ku_i$
$\partial_i=\frac{\partial}{\partial x_i}$
$\nabla=\partial_i$
$\nabla f=\partial_if$
$\nabla\cdot\vec F=\partial_iF_i$
$\nabla^2=\partial_i\partial_i$
$\nabla^2f=\partial_i\partial_if$
$\nabla^2\vec F=\partial_i\partial_i\vec F_j$
$(\nabla\times\vec F)_i=\varepsilon_{ijk}\partial_jF_k$
$\vec \nabla \times (\vec \nabla f) = \vec 0$
- Vector Calculus
- Jacobians and Change of Coordinate Systems
$$J=\varepsilon_{ijk}\partial iu\partial_jv\partial_kw=\frac{\partial(x_1,x_2,\ldots,x_n)}{\partial(a_1,a_2,\ldots,a_n}=\det\begin{bmatrix}
\frac{\partial x_1}{\partial a_1} & \cdots & \frac{\partial x_n}{\partial a_1}\\
\vdots & \ddots & \vdots\\
\frac{\partial x_1}{\partial a_n} & \cdots & \frac{\partial x_n}{\partial a_n}
\end{bmatrix}$$
$dV=|J|dx_1dx_2\ldots dx_n$
$J_{\text{polar}}=J_{\text{cylindrical}}=r$
$x=r\cos \theta, \space y=r\sin \theta, \space z=z$
$r=\sqrt{x^2+y^2},\space \theta=\tan^{-1}(\frac{y}{x}), \space z=z$
$ds^2=(dr)^2+(rd\phi)^2+(dz)^2$
$J_{\text{spherical}}=\rho^2\sin(\phi)$
$x=\rho\sin\phi\cos\theta,\space y=\rho\sin\phi\sin\theta,\space z=\rho\cos\phi$
$\rho=\sqrt{x^2+y^2+z^2},\space \theta=\tan^{-1}(\frac{y}{x}),\space \phi=\cos^{-1}(\frac{z}{\rho})$
$ds^2=(dr)^2+(r\sin\theta d\phi)^2+(rd\theta)^2$
- Line, Surface, Vector Integrals
- $\int_\Gamma\vec F\cdot d\vec l=\int_{t_0}^{t_f}\vec F(\vec r(t))r'(t)dt,\space r(t)=(y(t),x(t))\text{ where }y=f(x)\text{ describes }\Gamma$
Called a line integral (or vector line integral more specifically)
- $\int_S \vec F\cdot\hat ndA=\int_{b_0}^{b_f}\int_{a_0}^{a_f}(\vec F\cdot\hat n)dadb$
Surface Integral
- Major Theorems
$\int_\mathscr{C}\nabla f\cdot d\vec r=f(S_f)-f(S_0)\\\text{ (Fundamental Theorem of Vector Integrals)}$
$\iint_S(\vec F\cdot \hat n)dS=\iiint_V(\vec\nabla\cdot\vec F)dV\text{ (Divergence Theorem})$
$\iint_S(\vec\nabla\times\vec F)\cdot\hat ndA=\oint_\Gamma\vec F\cdot d\vec l\text{ (Stoke's Theorem})$
- Generalized Vector Spaces
$V \coloneqq \{u,v,\ldots\} :$ addition and multiplication by numbers are defined:
$\ket{u}, \ket{v} \in V \rightarrow \alpha \ket{u} + \beta \ket{v} \in V, \alpha,\beta \in \mathbb{R} \text{ or } \alpha,\beta \in \mathbb{C}$
Properties of vector spaces:
- (1) $\exists \ket{0} : \ket{u} + \ket{0} = \ket{u}$
There is a zero vector such that adding it to any other vector results in no change
- (2) $\exists \ket{-u} : \ket{u} + \ket{-u} = \ket{0}$
There is such thing as the negative of a vector, which when added to its positive counterpart, results in the zero vector
- (3) $\text{Vectors linearly independent} \leftrightarrow (\sum \alpha_i \ket{u_i} = \ket{0} \rightarrow \alpha_i = 0 ~ \forall i)$
A set of vectors are linearly independent if the only way for a linear combination of them to equal the zero vector is if their coefficients are all 0
i.e. Linearly independent vectors cannot form linear combinations of each other
- (4) $\text{Vectors } \{\ket{e_i}\} \text{ are a basis of }V \leftrightarrow \text{linearly independent and } \exists \alpha_i : \ket{u} = \sum_i \alpha_i \ket{e_i} \forall \ket{u} \in V$
A basis is a set of vectors that are each linearly independent and can form linear combinations to become any other vector in $V$
An example is $\{\sin(\omega t\}, \cos(\omega t)\}$ for the space of solutions of $\ddot x + \omega^2 x(t) = 0$
- (5) $\text{dim}(V) = |\ket{e_i}|$
The dimension of a vector space is the number of elements in any of its bases
- (6) $V_1 \cong V_2 \leftrightarrow (\exists L : V_1 \rightarrow V_2, \text{ an invertible linear operator})$
Two vectors spaces are isomorphic if there exists a linear operator $L$ between them, meaning there is a bijection between the two vector spaces and $L(\alpha \ket{u_1} + \beta \ket{v_1}) = \alpha L \ket{u_1} + \beta L \ket{v_1} = \alpha \ket{u_2} + \beta \ket{v_2} \forall \alpha,\beta$
Examples:
$\mathbb{R}^n$
$\mathbb{C}^n$
$\mathbb{P}^n = \{\text{polynomials of order} \le n\}$
$\{\text{solutions to linear homogeneous differential equations}\}$
- Infinite Dimensional Vector Space
E.g.:
Set of continuous functions $f : [0, 1] \rightarrow \mathbb{R}$
Set of square integrable functions on $[0, 1]$, or $\{f : \int_0^1 (f(x))^2 dx\}$
Set of all polynomials of degree $n$ where $n$ may be infinite
Projection Operator: $L^2 = L$
Scalar Product:
- $( ~ , ~) : V \times V \rightarrow \mathbb{C}$
$(a,b) = \braket{a|b}$
- (1) $(a,b) = (b,a)^*$
Scalar products are not quite commutative (have to take the complex conjugate)
- (2) $(a, \alpha b + \beta c) = \alpha(a,b) + \beta(a,c)$
Scalar products are distributive
- (3) $\text{For orthongonal vectors} \ket{u},\ket{v} : \braket{u | v} = 0$
Orthogonal vectors have a scalar product of 0
- (4) $\text{For } \{\ket{e_i}\}, \braket{e_i | e_j} = \delta_{ij}$
Basis vectors of the same index have a scalar product of 1, 0 if they have different indices
E.g.:
For square integrable functions on $[0,1]$
$(f,g) = \int_0^1 fg ~ dx$ on $\mathbb{R}$
Or
$(f,g) = \int_0^1 f^*g ~ dx$ on $\mathbb{C}$
E.g.:
$\{\frac{1}{\sqrt{2L}}, \frac{1}{\sqrt{L}}\sin(\frac{n\pi}{L}x), \frac{1}{\sqrt{L}}\cos(\frac{n\pi}{L}x)\} \text{ is an orthonormal basis}$
$(f,g) = \int_{-L}^Lfg ~ dx$
E.g.:
$\{\frac{1}{\sqrt{2\pi}} e^{inx}, n \in \mathbb{Z}\} \text{ is an orthonormal basis}$
$(f,g) = \int_{-\pi}^\pi f^*g ~ dx$
- $\braket{e_i | f} = \alpha_i$
$e_i$ is a basis vector
$f$ is a generalized vector (can be a function)
$\alpha_i$ is the coordinate of $f$ for the $i$th basis vector
- $(a,a) = \lVert a \rVert^2$
- Fourier Analysis
$f(x + \lambda) = f(x) ~ \forall x$
$\bar f = \frac{1}{b - a} \int_a^b f(x)dx$
$f(x) = \frac{1}{2}a_0 + \displaystyle\sum_{n = 1}^{\infty} a_n\cos(nx) + b_n\sin(nx)$
$a_0 = \frac{1}{\pi}\int_0^{2\pi}f(x)dx$
$a_n = \frac{1}{\pi}\int_0^{2\pi}f(x)\cos(nx)dx$
$b_n = \frac{1}{\pi}\int_0^{2\pi}f(x)\sin(nx)dx$
Generalized:
For $f(x) = f(x + 2L)$
$f(x) = \frac{a_0}{2} + \displaystyle\sum_{n = 1}^\infty\left[a_n\cos\left(\frac{n\pi}{L}x\right) + b_n\sin\left(\frac{n\pi}{L}\right)\right]$
$a_0 = \frac{1}{L}\int_0^{2L}f(x)dx$
$a_n = \frac{1}{L}\int_0^{2L} f(x) \cos(\frac{n\pi}{L}x)dx$
$b_n = \frac{1}{L}\int_0^{2L} f(x) \sin(\frac{n\pi}{L}x)dx$
Using $\omega$:
For $f(x) = f(x + \frac{2 \pi}{\omega})$
$f(t) = \frac{a_0}{2} + \displaystyle\sum_{n = 1}^\infty[a_n\cos(n \omega t) + b_n\sin(n \omega t)]$
$a_0 = \frac{\omega}{\pi}\int_0^{\frac{2\pi}{\omega}}f(t)dt$
$a_n = \frac{\omega}{\pi}\int_0^{\frac{2\pi}{\omega}} f(t) \cos(n \omega t)dt$
$b_n = \frac{\omega}{\pi}\int_0^{\frac{2\pi}{\omega}} f(t) \sin(n \omega t)dt$
$f(x) = \displaystyle\sum_{n = -\infty}^{\infty} C_ne^{inx}$
$f^*(x) = \displaystyle\sum_{n = -\infty}^{\infty} C_n^*e^{-inx}$
- $C_n = \frac{1}{\sqrt{2\pi}} \int_{-\pi}^{\pi}fe^{-inx}dx$
$C_n$ is like a coordinate of $f$ for the basis vector $e^{inx}$
$C_n = \braket{e_n | f} = \braket{\frac{1}{\sqrt{2\pi}} e^{inx} | f} = \frac{1}{\sqrt{2\pi}} \int_{-\pi}^{\pi}fe^{-inx}dx$
- Fourier Transform
$\hat f(k) = \int_{-\infty}^\infty f(x)e^{-2\pi ixk}dx$
$F(\omega) = \mathcal{F}(f(x)) = \frac{1}{2\pi}\int_{-\infty}^{\infty}f(x)e^{i\omega x}dx$
$f(x) = \frac{1}{\sqrt{2 \pi}}\int_{-\infty}^\infty a(k)e^{ikx}dk$
$a(k) = \frac{1}{\sqrt{2\pi}}\int_{-\infty}^\infty f(x) e^{-ikx}dx$
- Laplace Transform
$\hat x(s) = \int_0^\infty dt e^{-st}x(t)$
$T(x,t) = \int_{-\infty}^\infty G(x, x') f(x')dx'$
$G(x - x') = \frac{1}{\sqrt{4\pi Dt}}e^{-\frac{(x - x')^2}{4Dt}}$
- Convolutions
3B1B Video: https://www.youtube.com/watch?v=KuXjwB4LzSA
- $$(a * b)_n = \displaystyle\sum_{\substack{i, j \\
i + j = n}}
a_i \cdot b_j = \displaystyle\sum_{i = 1}^n a_i \cdot b_{n - i}$$
The convolution of $a_i$ and $b_i$ is given by the above formula
- Legendre Series
- Differential Equations
- $\nabla^2\phi = 0$
- $\frac{1}{c^2}\frac{\partial^2 f}{\partial t^2} = \nabla^2 f$
Called the wave equation
- $\frac{\partial T}{\partial t} = \alpha \nabla^2T$
Called the heat equation
3b1b video: https://www.youtube.com/watch?v=ToIXSwZ1pJU
$T$ is temperature and is a function of position and time (e.g.: $T(\vec x, t)$)
$\alpha$ is the thermal diffusivity of the material with units $\frac{m^2}{s}$
To solve for $T$
(1) Separation of variables (check differential equations section for more info)
(2) Apply boundary conditions
(3) Apply initial conditions
- Lagrange’s Method
- $J = \int_{x_1}^{x_2}f\{y(x), y'(x); x\}dx$
The goal is to find $y(x)$ to minimize $J$
- $\frac{\partial \mathcal{L}}{\partial q_i} - \frac{d}{dt}\frac{\partial L}{\partial \dot q_i} = 0, ~ i = 1, 2, 3$
Euler-Lagrange Equation
$\frac{\partial \mathcal{L}}{\partial q_i} = \frac{d}{dt}\frac{\partial \mathcal{L}}{\partial \dot q_i}$
$(\frac{\partial f}{\partial Y} - \frac{d}{dt}\frac{df}{d\dot Y})(\frac{\partial g}{\partial Y})^{-1} = (\frac{\partial f}{\partial z} - \frac{d}{dt}\frac{\partial f}{\partial \dot z})(\frac{\partial g}{\partial z})^{-1}$
- Hamiltonian
- $\delta \int_{t1}^{t2}L dt = 0$
Called Hamilton’s Principle, or Principle of Least Action
- $\displaystyle H \equiv \sum_i p_i \dot q_i - \mathcal{L}(q_i, \dot q_i, t), ~ p_i \equiv \frac{\partial \mathcal{L}}{\partial \dot q_i}, ~ \dot q_k = \frac{\partial H}{\partial p_k}, ~ \dot p_k = -\frac{\partial H}{\partial q_k}$
Hamilton’s equations of motion
- $\mathcal{L}(x, \lambda) \equiv f(x) + \lambda g(x)$
Lagrange undetermined multipliers
- Quantum Physics
- $\omega = 2\pi f = \frac{2\pi}{T}$
This is the angular temporarl frequency
- $k = 2\pi \xi = \frac{2\pi}{\lambda}$
This is the angular spatial frequency
$v_\text{phase} = \frac{\omega}{k} = f \lambda = \frac{\lambda}{\tau} = \frac{2\pi f}{k} = \frac{\omega \lambda}{2 \pi}$
$v_\text{group} = \frac{\partial E}{\partial p} = \frac{\partial \omega}{\partial k}$
$v_\text{phase}v_\text{group} = c^2$
$v = \frac{\omega}{|k|\cos\theta} = \frac{\lambda}{T\cos\theta}$
$E = \hbar\omega = \frac{h}{T} = hf$
$p = \hbar k = \frac{h}{\lambda} = h\xi$

- $n\lambda = d\sin \theta$
The equation for Bragg scattering
$n$ is the diffraction order
$\lambda$ is the electron wavelength
$d$ is the lattice period
$\theta$ is the scattering angle
- $\lambda = \frac{\lambda_C}{\sqrt{\frac{2K}{mc} + (\frac{K}{mc^2})^2}}, ~ \lambda_C = \frac{h}{mc} = \frac{2\pi}{k_C}$
$\lambda_C$ is the Compton wavelength of a particle
$\lambda_{C_\text{electron}} = 2.43 pm$
The compton wavelength of a certain particle measures the wavelength of a photon that matches the energy of that certain particle’s rest energy
$m_ec^2 \approx 511 keV$
- $\frac{1}{\lambda_t^2} - \frac{1}{\lambda^2} = \frac{1}{\lambda_C^2}$
$\lambda_t$ is the temporal wavelength
$\lambda$ is the de Broglie wavelength
$\lambda_C$ is the Compton wavelength of the electron
- $R = \frac{\lambda}{2\text{NA}}$
$R$ is the spatial resolution of the microscope
$\text{NA}$ is the numerical aperature
$\frac{e^2}{4\pi\varepsilon_0} = 1.44 ~ eV \cdot nm$
$\hbar c = 197 ~ eV \cdot nm$
$\alpha = \frac{e^2}{4\pi\varepsilon_0\hbar c} \approx \frac{1}{137}$
$E\lambda \approx 400\pi ~ eV\cdot nm$
- $\Delta t \Delta \omega \ge \frac{1}{2}, ~ \Delta x \Delta k \ge \frac{1}{2}$
Unitless
- $\Delta t \Delta E \ge \frac{\hbar}{2}, ~ \Delta x \Delta p \ge \frac{\hbar}{2}$
Heisenberg Uncertainty Principle
Units of action
$f(x) = \frac{1}{\sqrt{2\pi}}\int_{-\infty}^\infty a(k)e^{ikx}dk$
$a(k) = \frac{1}{\sqrt{2\pi}}\int_{-\infty}^\infty f(x)e^{-ikx}dx$
$\text{To time: }x\rightarrow t, ~ k\rightarrow \omega, ~ i \rightarrow -i$
$\frac{}{}$$f(x) = \frac{1}{2\pi}\iint f(x')e^{-ikx'}dx'e^{ikx}dk = \frac{1}{2\pi}\iint f(x')e^{ik(x - x')}dkdx' = \int f(x')\delta(x - x')dx' = f(x)$
$\delta(x - x') = \frac{1}{2\pi}\int_{-\infty}^\infty e^{ik(x - x')}dk$
$f(x) = \int_a^b f(x')\delta(x - x')dx'$
- $-\frac{1}{c^2}\frac{\partial^2 u}{\partial t^2} + \nabla^2 u = k_c^2u$
Wave equation
$\Psi(x,0) = \frac{1}{\sqrt{2\pi}}\int_{-\infty}^\infty \Psi(k)e^{ikx}dk$
$\Psi(x,t) = \frac{1}{\sqrt{2\pi}}\int_{-\infty}^\infty \Psi(k)e^{i(kx-\omega(k)t)}dk, ~ \omega = \frac{\hbar k^2}{2m}$
$-\frac{\hbar^2}{2m}\frac{d^2 \Psi}{dx^2} + U(x)\Psi(x, t) = i\hbar \frac{\partial \Psi}{\partial t}$
$-\frac{\hbar^2}{2m}\frac{d^2 \Psi}{dx^2} + U(x)\Psi(x) = E\Psi(x)$
$E_n = \frac{\hbar^2 k^2}{2m} = \frac{n^2 \pi^2 \hbar^2}{2mL^2}, ~ n \in \mathbb{Z}^+$
$\Psi_n(x) = A\sin(\frac{n\pi x}{L})$
$\Psi_0(x) = \frac{1}{\sqrt{\sigma_x \sqrt{2\pi}}}e^{-(\frac{x}{2\sigma_X})^2} = \left(\frac{m\omega_0}{\pi\hbar}\right)^{\frac{1}{4}}e^{-\frac{m\omega_0 x^2}{2\hbar}}$
$\Psi_0(k) = \frac{1}{\sqrt{\sigma_x \sqrt{2\pi}}}e^{-(\frac{k}{2\sigma_X})^2} = \left(\frac{\hbar}{\pi m \omega_0}\right)^{\frac{1}{4}}e^{-\frac{\hbar k^2}{2m\omega_0}}$
$P(x \in [a, b]) = \int_a^b |\psi(x)|^2dx$
- $\braket{x} = \int_{-\infty}^\infty x|\Psi (x, t)|^2 dx$
$\braket{x}$ is the expected value of the position of a particle
$|\Psi (x, t)|^2$ is the same as $P(x)$, the probability density of the particle
$\sigma_x^2 = (\Delta x)^2 = \braket{x^2} - \braket{x}^2$
$E_n = (n + \frac{1}{2})\hbar \omega$
- $\Psi(x, t) = Ae^{i(kx - \omega t)} + Be^{i(-kx - \omega t)}$
This wave function describes the wave traveling toward or away from a potential energy barrier.

If the wave source is on the left side of the barrier:
For the $\Psi$ where $x < 0$:
$A, B \in \mathbb{C}$
For the $\Psi$ where $x > 0$:
$B = 0$ because $e^{i(-kx - \omega t)}$ for $k > 0$ represents a wave traveling leftward, which shouldn’t be possible: https://www.desmos.com/calculator/pew0elwjfz
- $R = \frac{(\Psi * \Psi)_\text{reflected}}{(\Psi * \Psi)_\text{incident}} = \frac{|B|^2}{|A|^2}$
The reflection coefficient $R$ is the ratio of the reflected probability density to the incident probability density
- $T = \frac{(\Psi * \Psi)_\text{transmitted}}{(\Psi * \Psi)_\text{incident}} \cdot \frac{k_2}{k_1} = \frac{|F|^2}{|A|^2}$
The transmission coefficient $T$ is the ratio of the transmitted probability density to the incident probability density
- $R + T = 1$
A wave is either reflected or transmitted, so the reflected and transmitted probabilities must add up to $1$
- $k = \frac{\sqrt{2m(E - U_0)}}{\hbar}$
The wave number for a particle moving through space
- $\lambda_f - \lambda_i = \Delta \lambda = \frac{h}{m_e c}(1 - \cos\theta)$
Compton scattering equation
$\lambda_f$ is the wavelength of the scattered photon
$\lambda_i$ is the wavelength of the incident photon

$m_e vr = n\hbar$
$2\pi r = n\lambda$
- $r_n = \frac{n^2\hbar^2}{m_ek e^2} = n^2a_0$
$r_n$ is the radius for the orbit of an electron at energy level $n$ orbiting a nucleus
$n$ is the energy level, $\in \mathbb{Z}^+$
$k$ is Coloumb’s constant
$e$ is the charge of the electron
- $a_0 = \frac{\hbar^2}{m_e k e^2} = \frac{1}{\alpha k_e} = \frac{\lambda_e}{2\pi\alpha} \approx 0.529 Å$
$a_0$ is the Bohr radius for when $n = 1$
This is the smallest radius for an electron to be orbiting a hydrogen atom
- $E_i = \frac{\vec p_i^2}{2m} + \frac{1}{2}k\vec r_i^2$
This is the equipotential theorem (might be spelling it wrong)
- $2\braket{K} = n\braket{U}$
Virial Theorem
- $E_n = -\frac{ke^2}{2a_0}(\frac{1}{n^2})$
The energy level for an electron at energy level $n$
$r_n = n^2\frac{a_0}{Z}$
$E_n = -\frac{ke^2}{2a_0}(\frac{Z^2}{n^2})$
$\mu = \frac{m_e m_p}{m_e + m_p}$
$f_\text{MB} = Ae^{-\frac{E_i}{k_BT}}$
- $K_\text{max} = hf - \phi$
Photoelectric effect
States that the kinetic energy of a photoelectron is dependent on the frequency of the light hitting the material and work function of the material
$\phi$ is the work function
- $j^\star = \sigma T^4$
Stefan’s Law
$j^\star$ is black-body radiant emittance in terms of intensity, or $\frac{W}{m^2}$
$\sigma$ is the Stefan-Boltzmann constant and has a value of $\frac{2\pi^5 k^4}{15 c^2 h^3} = 5.670374419 * 10^{-8} ~ W m^{-2} K ^{-4}$
$T$ is thermodynamic temperature
$u(\lambda, T) = \frac{8 \pi hc}{\lambda^5 (e^{hc / \lambda k T} - 1)}$
- Particle Physics
- Standard Model
Video explaining: https://www.youtube.com/watch?v=mYcLuWHzfmE
Infographic: https://www.flickr.com/photos/95869671@N08/51148317732/

| Particle | Fermion | Boson |
| | | |
| Mass | Has mass | No mass |
| Spin | $\frac{1}{2} + k, k \in \mathbb{N}^+ \\
\frac{1}{2}, \frac{3}{2}, \frac{5}{2}, \ldots$ | $k, k \in \mathbb{Z} - \mathbb{Z}^- \\
0, 1, 2, \ldots$ |
- $\text{Baryon number} = \frac{1}{3}(n_{q} + n_q')$
$n_q$ is the number of quarks in the baryon
$n_q'$ is the number of anti quarks in the baryon
A baryon is made of an odd number of 3 or more quarks
- $\text{Color charge}$

There are 3 color charges that quarks and gluons can take: red, green, and blue
There are 3 color charges that antiquarks and antigluons can take: anti-red (cyan), anti-green (magenta), and anti-blue (yellow)
For protons/neutrons, you need a neutral overall color charge, meaning 3 or each color for each quark
For pions, you also need a neutral overall charge, meaning 1 quark and 1 anti quark of opposing colors
- Thermodynamics
$k_B = 1.380649*10^{-23}\frac{J}{K}$
- Laws
0) If 2 systems are in thermal equilibrium with a 3rd, they are all in the thermal equibilbrium with each other (transitivity, thermometers exists)
1) Energy is conserved
2) In process, the total entropy change $\Delta S > 0$
3) The entropy $S \rightarrow k$ a constant (usually $0$) as $T \rightarrow 0$
Analogy to life:
0: There is a game
1: You can’t win
2: You can’t break even
3: You have to play
$\Omega(E) =$ # of distinguishable states accessible w/ energy E in $[E, E + dE]$
$S = k_b \ln(\Omega (E))$
$\frac{1}{T} = \frac{\partial S}{\partial E}$
$N_{MB} = \frac{N!}{n_1! n_2! n_3! \ldots}$, $n_i$ are the indistinguishable outcomes, $\sum^+ n_i = N$
$\Omega(E) = \sum_\text{arrangements}N_{MB}$
$\bar n_j = n_{j1}p_1 + n_{j2}p_2 + \ldots + n_{jm}p_m$
$\Omega = {\text{\# energy states + \# particles - 1} \choose \text{\# energy states}}$
$PV = Nk_BT$
$K_\text{gas} = \frac{3}{2}k_BT$
$n(v)dv = 4\pi \frac{N}{V}(\frac{m}{2\pi k_BT})^{3/2}v^2e^{-mv^2/2k_BT}dv$
$v_\text{mp} = \sqrt{\frac{2k_BT}{m}}$
Units μ Molarity=Lmol 1ppm=1Lmg w/v%=100mL of solutiongrams of material
Uncertainty ± C=A±B→CΔC=CΔA+CΔB C=AB or C=BA or C=AB→CΔC=AΔA+BΔB
Critical Ratio Test
A test to determine whether to eliminate a questionable result as an outlier for a small number of trials (n<7) Rq=xq−xfxq−xn Rq is the deviation ratio xq is questionable result xn is the nearest result to the questionable result xf is the furthest result from the questionable result
If Rq exceeds the critical ratio designated by the chart below, then the questionable result may be excluded as an outlier
RAD=100%xˉΔx,Δx=n∑i=1n∣(xi−xˉ)∣
This is the relative average deviation (RAD) of a sample xi is a sample point’s value xˉ is the average of the sample n is the number of samples
E.g.: the RAD of {1,2,3,4,5} is 100%⋅5∣5−3∣+∣4−3∣+∣3−3∣+∣2−3∣+∣1−3∣=120%
Constants R∞ u=1.6605390660⋅10−24g (Atomic mass unit) NA=6.022142⋅1023molparticles (Avogadro’s number) kB=1.380651⋅10−23KJ (Boltzmann’s constant) R=0.0820575mol⋅KL⋅atm=8.31447mol⋅KJ (Gas constant) e=1.60217634⋅10−19C (Charge on electron) F=9.6485538⋅104mole−C (Faraday’s constant) me=9.109383⋅10−28g (Mass of electron) mn=1.6726217⋅10−24g (Mass of neutron) mp=1.6726217⋅10−24g (Mass of proton) π=3.14159265358979323846… h=6.626069⋅10−34J⋅s (Planck’s constant) ℏ=2πh=1.05457266⋅10−34J⋅s (Reduced Planck’s constant) c=2.99792458⋅108sm (Speed of light in vacuum, exact) ε0=4πk1=8.85418782⋅10−12N⋅m2C2(mF) μ0=4π⋅10−12A2N
Bohr Model ⚛ Zeff≈Z−inner shell e−−21 number of same shell e− rn=Zeffn2a0,n∈N
a0=e2πmeh2ε0=0.529177A˙
Called a Bohr radius, or the radius at which an electron would most likely to be found in a hydrogen atom A˙=10−10m e=1.60217663∗10−19C me=9.109∗10−31kg En=−R0n2Zeff2
σxσp≥4πh
This is called Heisenberg’s Uncertainty Principle σx is the uncertainty in position, σp is the uncertainty in momentum
For more massive objects, σx is very small σp=mσv σ=a0Zr
H^Ψ(x)=EΨ(x)=−8π2mh2dx2d2Ψ(x)+U(x)Ψ(x)
The kinetic energy term is 8π2mh2dx2d2Ψ(x) while the potential energy term is U(x)Ψ(x)
De Broglie’s Principle: λ=ph λ is the object’s associated wavelength h=6.626069⋅10−34J⋅s (Planck’s constant) p is the momentum of the object
So for macroscopic objects, p is huge, so λ shrinks to negligible values
For nanoscopic objects, p is smaller, so λ is no longer negligible P(x)=∣Ψ(x)∣2 ∫−∞∞∣Ψ(x)∣2dx=1 P(r)=4πr2(Rn,o(r))2dr≈Ψ2dV
1D PIB
For a 1D Particle in a Box: Ψ(x)=L2sin(Lnπx),n∈N En=8mL2n2h2,n∈N
L2=ℓ(ℓ+1)4π2h2,ℓ=0,1,2,…,n−1 L is the angular momentum while ℓ is the angular momentum quantum number ℓ=0→s,1→p,2→d,3→f Lz=2πmlh=mlℏ,ml=0,±1,±2,±ℓ ml is called the magnetic quantum number rnlˉ=Zeffn2a0{1+21[1−n2ℓ(ℓ+1)]} Zeffns>Zeffno>Zeffnd>Zeffnf
(n,l,ml,ms) n∈N is the principle quantum number and it dictates the orbital size and energy level l∈N:[0,n−1] is the secondary (angular momentum) quantum number and it dictates the subshell type (s, p, d, or f) ml∈N:[−l,l] is the magnetic quantum number and it dictates the orbital orientation (x, y, z, or etc.) ms∈{−21,21} is the spin quantum number and it dictates the electron spin direction
\mu=|q|r\space(C\cdot m)$$\mu=qr\space(C\cdot m)
The dipole moment between two charges is equal to the separated charge times the distance betwen them (e.g. bond distance)
Spectroscopy 🌈 ΔE=hν c=λν
Photon Type:
Radio
Microwaves
Infrared
Visible
UV
X-rays
Gamma
λmin
1m
1mm
700nm
400nm
10nm
10pm
N/A
Emin(cm−1)
0.001
1
1000
16000
100000
1E+8
8E+9
Emin(molkJ)
0.0030125
0.0125
12.5
200
1250
1250000
1E+8
Bonds
Nuclear spin
Rotational
Vibrational
Valence
Valence
Inner shell
Nuclear
- $E_\text{rotational}(J)=BhJ(J+1),\space J\in\mathbb{N}$
A permanent dipole moment is required for rotational energy states to be available because otherwise, the oscillating $\vec E$ from the photon of light will not be able to excite the molecules to a higher rotational state.
- $\Delta E=2BhJ_\text{larger},\space J\in\mathbb{N}$

- $B=\frac{h}{8\pi^2I}$
Units are $Hz$ because $h$ is the $J\cdot s$, or $kg\frac{m^2}{s^2}\cdot s=kg\frac{m^2}{s}$, and when divided by $I$, which is in units of $kg\cdot m^2$, $B$ becomes $\frac{1}{s}$, or $Hz$
- $\tilde B=\frac{B}{c}$
$\tilde B$ is the wavenumber form of $B$, and $c$ is $2.998\cdot 10^{10}\frac{cm}{s}$.
Units go from $Hz$ to $\frac{1}{cm}$
In physics, the SI units for wavenumber are $\frac{1}{m}$ rather than $\frac{1}{cm}$ in chemistry. It is also calculated as $K=\frac{2\pi}{\lambda}$ in physics rather than $\bar v=\frac{1}{\lambda}$ in chemistry.
$I=\mu R^2$
$\mu=\frac{m_1m_2}{m_1+m_2}$
- $E_\text{vibrational}(V)=(\frac{1}{2}+V)h\nu,\space V\in\mathbb{N}$
The molecule needs to either be very massive (e.g. $I_2$) or the temperature needs to be very high for vibrational states to be excited.
$\Delta E=h\nu$
$\nu=\frac{1}{2\pi}\sqrt{\frac{k}{\mu}}$
$\tilde\nu=\frac{\nu}{c}$
- $A=\varepsilon lc=-\log(T),\space T=\frac{I}{I_0}$
Beer-Lambert Law
$A$ is the absorbance
$\varepsilon$ is the molar absorptivity and is constant only for values of $A$ between 0.1 and 1.0
$l$ is the length the light has to pass through for the solution
$c$ is the concentration of the solution
$T$ is the transmittance
- Gases 💨
- $P_\text{total}=P_a+P_b=X_a{}^{0}P_a+X_b{}^{0}P_b$
This is called Raoult’s Law
$X$ is the mole fraction of a gas
${}^{0}P$ is the vapor pressure of a gas
- $P_a=kX_a$
This is called Henry’s Law
The interpretation is that the higher the partial pressure $P_a$ of a gas over a liquid, the higher molefraction $X_a$ of the gas dissolved in the liquid.
Only really works for very dilute/small $X_a$, then it follows something like Raoult’s Law
$k$ is some constant $>{}^{0}P_a$
- $P_\text{vapor}\propto T_\text{boiling}$
Vapor pressure and boiling point are inversely related
$\text{Boyle's Law: }V\propto \frac{1}{P}$
$\text{Charle's Law: }V\propto T$
$\text{Moles: }V\propto n$
- $PV=nRT$
$R=0.082057\frac{L\cdot atm}{K\cdot mol}=8.314472\frac{J}{K\cdot mol}$
$\frac{101.33J}{L\cdot atm}$
Assumptions (valid when $T$ is high and $P$ is low):
Particle volume is negligible
The particles follow straight line paths until collision
$\Delta K=0$ (perfectly elastic collisions)
There are no intermolecular forces
$v_{\text{most probable}}=\sqrt{\frac{2k_bT}{m}},\space \bar v=\sqrt{\frac{8k_bT}{\pi m}},\space v_{\text{root mean square}}=\sqrt{\frac{3k_bT}{m}}$
- $(P+a\frac{n^2}{V^2})(V-nb)=nRT$
$a$ is the attractive force, $b$ is the repulsive force
$d=\frac{P\cdot MM}{RT}\space(g/L)$
- $K_\text{gas}=\frac{3}{2}nRT=\frac{3}{2}Nk_BT$
$n$ is the number of moles of gas
$N$ is the number of particles of gas
$K_\text{particle}=\frac{3}{2}k_BT$
- $Z=\frac{PV}{nRT}$
$Z$ approaches 1 as the gas behaves more like an ideal gas
- $\%\text{ composition of gas}_i = (100\%) (\frac{A_i}{A})$$
For gas chromatography, the percent composition a species $i$ is equal to the fraction of its area of the chromatogram peak over the combined area of all peaks
Area is calculated as the width at half the height of the peak times the height of the peak
Hydrogen Bonding: NOF
- Thermochemistry Δ
$\Delta H=\Delta U+\Delta(PV)$
$U_\text{tot}=U_\text{trans}+U_\text{rot}+U_\text{vib}+U_\text{elec}$
$\text{df}_\text{total}=3N$
$\text{df}_\text{trans}=3$
$\text{df}_\text{rot}=3\text{ (or 2 if linear)}$
$\text{df}_\text{vib}=3N-6\text{ (or } 3N-5 \text{ if linear)}$
$U_\text{trans}=\frac{\text{df}_\text{trans}}{2}RT$
$U_\text{rot}=\frac{\text{df}_\text{rot}}{2}RT$
$U_\text{vib}=\text{df}_\text{vib}RT\text{ (only for massive molecules or high T)}$
- $\Delta U_\text{system}=q+W_\text{by system}$
Isochoric: $\Delta V=W=0$
Isobaric: $\Delta P=0$
Adiabatic**: $\Delta q=0,\space \Delta U=W$**
Isothermal: $\Delta T=\Delta U=0$
- $W_\text{by system}=-W_\text{by environment}=-\int_{V_0}^{V_f}PdV$
This is the opposite of in physics, as $W_\text{by system}=\Delta U$ in chemistry rather than $W_\text{by environment}=-\Delta U$ in physics
Units of $L\cdot atm$ are inconvenient and can be translated by $\frac{101.325J}{L\cdot atm}$
- $W_\text{by irreversible process}=-P_\text{ext}\Delta V$
$P_\text{ext}$ is the external pressure and also the final pressure of the system
$W_\text{by isochoric system}=0$
$W_\text{by isobaric system}=-P\Delta V$
$W_\text{by adiabatic system}=nc_v\Delta T=-P_\text{ext}\Delta V$
$W_\text{by isothermal system}=-nRT\ln\left(\frac{V_f}{V_0}\right)$
$T_fV_f^{\gamma-1}=T_0V_0^{\gamma -1},\space\gamma=\frac{C_v}{C_p}\text{ for adiabatic systems}$
$P_fV_f^{\gamma-1}=P_0V_0^{\gamma -1},\space\gamma=\frac{C_v}{C_p}\text{ for adiabatic systems}$
$C=\frac{q}{\Delta T},\space c_n=\frac{q}{\Delta T\cdot n},\space c_s=\frac{q}{\Delta T\cdot g}(\text{same as }q=mc\Delta T)$
$C_v=\frac{dU}{dT},\space C_p=\frac{dH}{dT}=C_v+R$
Hess’s Law: $\Delta H_\text{net}=\Delta H_1+\Delta H_2+\ldots\Delta H_n$
- $\Delta H=\displaystyle\sum_{i=0}^{N}(H_\text{bonds broken}-H_\text{bonds formed})$
Breaking bonds requires energy and results in a positive contribution to enthalpy
Forming bonds releases energy and results in a negative contribution to enthalpy
- $\Delta H_\text{sublimation}\approx\Delta H_\text{fusion}+\Delta H_\text{vaporization}$
$\Delta H_\text{fusion}$ is the heat required to transition from solid to liquid
$\Delta H_\text{vaporization}$ is the heat required to transition from liquid to gas
$\Delta H_\text{sublimation}$ is the heat required to transition from solid to gas
- Thermodynamics 🔥
$\text{1st Law: }\Delta U_\text{uni}=\Delta U_\text{sys}+\Delta U_\text{sur}=0$
$\text{2nd Law: }\Delta S_\text{sys}+\Delta S_\text{sur}=\Delta S_\text{uni}\ge0$
$\text{3rd Law: }S_{\text{pure substances at 0}\degree K}=0$
$\Delta S=\int_{q_0}^{q_f}\frac{dq_\text{rev}}{T}=\frac{q_\text{rev}}{T}$
$\Delta S_\text{isothermal}=-nRT\ln(\frac{P_f}{P_0})$
$\Delta S_\text{adiabatic}=0$
$\Delta S_\text{phase change}=\frac{q_\text{rev}}{T}=\frac{\Delta H}{T}$
$\Delta S_\text{v}=nc_v\ln(\frac{T_f}{T_0})$
$\Delta S_\text{p}=nc_p\ln(\frac{T_f}{T_0})$
$S=k_B\ln(\Omega)$
$\Omega\propto V^N\text{ for ideal gases}$
- $\Delta G=\Delta H-T\Delta S$
Called Gibb’s Free Energy
Can be derived from the 2nd law of thermodynamics

- Equilibrium ⚖️
- $\Delta G^\degree=-RT\ln(K)$
At $25\degree C$
$K=e^{-\frac{\Delta G^\degree}{RT}}$
- $\Delta G=\Delta G^\degree+RT\ln(Q)$
Called the Nernst Equation
$K_f=K_0e^{-\frac{\Delta H}{R}(\frac{1}{T_f}-\frac{1}{T_0})}$
- $K=\frac{\displaystyle\prod_{i=0}^{N_p}[P_i]^{p_i}}{\displaystyle\prod_{j=0}^{N_R}[R_j]^{r_j}}=\frac{[C]^c[D]^d}{[A]^a[B]^b}\text{ for }(\ce{aA +bB<=>cC +dD})$
$p$ is the coeffecient of some $P$ product species
$r$ is the coeffecient of some $R$ reactant species
- $K_p=K_c[RT]^{\sum_{i=0}^{N_P}(p_i)-\sum_{j=0}^{N_R}(r_j)}=K_c[RT]^{c+d-a-b}\text{ for }(\ce{aA +bB<=>cC +dD})$
$k_c\text{ is for concentration}$
$k_p\text{ is for pressure}$
- $\text{5\% rule}$
If we arrive at some form of $k_{eq}=\frac{x^2}{a\pm x}$, then $x=\sqrt{k_{eq}*a}$ if $x<0.05*a$
$K_\text{sp}=[B]^b[C]^c\text{ for }(\ce{aA_{(aq)}<=>bB_{(aq)} +cC_{(aq)}})$
- ICE Tables
E.g. $\ce{HA_{(aq)} + H2O_{(l)} <=> A-_{aq} + H3O+_{(aq)}, K_c = 1.2*10^{-3}}$
| Reaction | $\ce{HA}$ | $\ce{H2O}$ | $\ce{A-}$ | $\ce{H3O+}$ |
| | | | | |
| Initial | $0.5$ | $\varnothing$ | $0$ | $0$ |
| Change | $-x$ | $\varnothing$ | $+x$ | $+x$ |
| Equilibirium | $0.5-x$ | $\varnothing$ | $x$ | $x$ |
$K=\frac{[A^-][H_3O^+]}{[HA]}$
- Acid-Base 🍋
- $\text{Equivalent Weight}$
Equivalent Weight of Arrhenius Acids:
$= \frac{g/mol\text{ acid}}{\text{equivalents of} ~ H^+\text{disassociated per mol of acid}}$
Equivalent Weight of Arrhenius Bases:
$= \frac{g/mol\text{ base}}{\text{equivalents of} ~ OH^- \text{ disassociated per mol of base}}$
$K_w=K_aK_b=[H^+][OH^-]=1.008\cdot10^{-14}$
$pX=-\log_{10}(X)$
$\text{Strong Acids: }\ce{HCl, HBr, HI, HNO_3, HClO_3, HClO_4, H_2SO_4}$$\text{Strong Bases: }\ce{LiOH, NaOH, KOH, RbOH, CsOH, Ca(OH)_2, Sr(OH)_2,\\Ba(OH)_2,XO^{2-},XNH^{2-},S^{2-}}$
$\ce{[OH^-] + [H_3O^+] << [A^-] +[HA]}\text{ for buffer solutions to exist}$
- $pH=pK_a+\log\frac{[A^-]}{[HA]}$
This is called the Henderson Hasselbach equation
$pOH=pK_b+\log\frac{[HB]}{[B^-]}$
- Electrochemistry 🔋
Electrons travel from anode to cathode (AC)
AOCR (Anode has oxidation, cathode has reduction)
AnOreXic anode oxidation, CanceRous cathode reduction
Anions migrate to anodes, cations migrate to cathodes
- $E_\text{cell}^\degree=E_\text{cathode}^\degree-E_\text{anode}^\degree$
Units are $V$
- $\Delta G^\degree=-nFE^\degree_\text{cell}$
$n$ is least common multiple of the number of electrons exchanged
$F=9.6485538\cdot10^4\frac{C}{mol}\text{ (Faraday's constant)}$
- $E=E^\degree-\frac{RT}{nF}\ln(Q)$
$\frac{RT}{F}\approx0.0592$
Called the Nernst equation
!https://mmsphyschem.com/wp-content/uploads/2021/12/standard-reduction-potential-table_193-3076027.jpg
The X agent includes the entire chemical formula, not the particular atom that gets X’ed, where X refers to oxidation or reduction.
- Kinetics 🚄
$\text{rate} = -\frac{1}{a}\frac{d}{dt}[A] = -\frac{1}{b}\frac{d}{dt}[B] = \frac{1}{c}\frac{d}{dt}[C] = \frac{1}{d}\frac{d}{dt}[D] \text{ for }\ce{aA + bB -> cC + dD}$
$\text{rate}=k[A]^a[B]^b\text{ for }\ce{aA + bB <=> \ldots}$
| Order | Rate law | Integrated rate law | Half-life | Unit of rate constant | Graph |
| | | | | | |
| $0$ | $\text{Rate}=k[A]^0$ | $[A]_t=-kt[A_0]$ | $t_{1/2}=\frac{[A]_0}{2k}$ | $mol\space L^{-1}s^{-1}$ | $\text{[A] vs t; slope = -k}$ |
| $1$ | $\text{Rate}=k[A]^1$ | $\ln[A]_t=-kt+\ln[A]_0$ | $t_{1/2}=\frac{\ln(2)}{k}$ | $s^{-1}$ | $\text{ln[A] vs t; slope = -k}$ |
| $2$ | $\text{Rate}=k[A]^2$ | $\frac{1}{[A]_t}=kt+\frac{1}{[A]_0}$ | $t_{1/2}=\frac{1}{k[A]_0}$ | $L\space mol^{-1}s^{-1}$ | $\frac{1}{[A]}\text{ vs t; slope = k}$ |
| n | $\text{Rate}=k[A]^n$ | $(n-1)kt=\frac{1}{[A]^{n-1}}-\frac{1}{[A_0]^{n-1}}$ | $t_{1/2}=\frac{2^{n-1}-1}{k(n-1)[A]_0^{n-1}}$ | $(mol\space L^{-1})^{1-n}s^{-1}$ | $\frac{1}{[A]^{n-1}}\text{ vs t; slope = k}$ |
- $\ln(\frac{k_t}{k_0})=\frac{E_a}{R}(\frac{1}{T_0}-\frac{1}{T_t}),\space k_t=k_0e^{\frac{E_a}{R}(\frac{1}{T_0}-\frac{1}{T_t})}$
This is called the Arrhenius equation, and it relates rate constants at varying temperatures given the activation energy of the reaction
$E_a$ is the activation energy of the reaction
$R=8.31447\frac{J}{mol\cdot K}\text{ (Gas constant)}$
- Nuclear Chemistry ☢️
Nuclear reactions follow the rules of first order reactions
- Solid State Chemistry 🖥
- Organic Chemistry 💊
Quizlet of Common Compounds:
https://quizlet.com/432085522/organic-compounds-structure-nomenclature-flash-cards/
Compound Drawer:
https://chem-space.com/search
- $\text{Delocalization}$
Electrons can be “delocalized” or shared between more than two atoms
$\pi$ bonds with lone pairs of $e^-$ or $\pi$ bonds that aren’t separated by an atom (so can’t have two single bonds between two $\pi$ bonds, only at most 1) are conjugated, or allow for delocalization
A $p$ orbital is required to be used, meaning the atom is either $sp^2$ or rarely $sp$ hybridized
Electron pushing is prioritized by the steps:
1) negative charge
2) lone pairs
3) double/triple bonds
The resonance structures that are most common prioritize these criteria:
1) Filled valence shells
2) Maximum number of covlanet bonds
3) Least separation of unlike charges
4) Negative charges on more electronegative atoms
- $\text{Functional Groups}$

- $\Delta G^\degree=-RT\ln(K_\text{eq})$
$R=1.987\frac{cal}{mol\cdot K}=1.987\cdot 10^{-3}\frac{kcal}{mol\cdot K}=8.314\frac{J}{mol\cdot K}$
- $\text{Stereochemistry}$

- $\text{Optical Acitivty}$
- $[\alpha]_\lambda^T=\frac{\text{Observed rotation }(\degree)}{\text{Length (dm) }\times\text{ Concentration}}$
$[\alpha]_\lambda^T$ is the specific rotation of a substance in a polarimeter
The length is typically 1 dm, the concentration is typically 1 g/mL, and the wavelength $\lambda$ is typically 589 nm from the Sodium D line.
- $\text{Percent optical purity}=\frac{[\alpha]_\text{sample}}{[\alpha]_\text{pure enantiomer}}\times 100\%$
- $\text{Acid Base}$
$K_a=\frac{[H^+][A^-]}{[HA]}$
$\text{p}K_\text{x}=-\log(K_\text{x})$
- $\text{p}K_{\text{a}\ce{(H2O)}}\approx 16$

- $\text{p}K_\text{eq}=\text{p}K_\text{HA}-\text{p}K_{\text{BH}^+}$
If the pKa of an acid in the reactants side is lower than the pKa of the conjugate acid (of the base in the reactants side), then the reaction favors the products.
Low pKeq means high Keq
- $\text{p}K_a\text{ table}$

- $\text{Mechanisms}$
Alkenes
- Rearrangment (1,2 shifts)

- Hydrohalogenation

In the addition of a hydrogen halide to an alkene reaction (hydrohalogenation, hydrochloration, hydrobromation, hydroiodation), the halide adds to the most substituted carbon and the hydrogen binds to the other carbon with rearrangment

- Hydration

In the addition of water to alkene with acid catalyst reaction (hydration), the hydroxyl group adds to the most substituted carbon with rearrangment

- Halogenation

In the addition of halogen gas to an alkene (halogenation) along with carbon tetrachloride catalyst, one of the halogen atoms forms a ring with the alkene. The other negative halide then attacks from the opposite side of the halonium ion in what’s called an anti co-planar attack without rearrangment

- Halohydrin Formation

In the addition of halogen and water to an alkene, the end result is that a hydroxyl group is added to the more substituted carbon while the halogen is added in anti to the less substituted carbon

- Oxymercuration

In the addition of mercuric acetate to an alkene (oxymercuration), the net result is that an OH is added to the most substituted carbon while an H is added in anti without rearrangement (the addition of H and OH may be syn or anti)

- Hydroboration

In the addition of borane to alkene with the presence of hydrogen peroxide and sodium hydroxide (hydroboration), a hydroxyl is added to the lesser substituted carbon along with a hydrogen in syn addition


- Halohydrin Formation

In the addition of a halogen and water to an alkene (halohydrin formation), one of the halogen atoms adds to the less substituted carbon while the hydroxyl adds to the more substituted carbon in an anti-coplanar attack without rearrangement

- Diol formation (oxidation)

In the addition of osmium tetraoxide to alkene (diol formation) along with sodium hydrogen sulfite, the end result is that two hydroxyl groups are added in syn addition to the previously double-bonded carbons to form a glycol (or vicinal diol).
- Ozonolysis (oxidation)

When ozone and DMS are reacted with alkene (ozonolysis), the carbon double bond is cleaved and an oxygen is added to both ends to form a ketone or aldehyde

- Hydrogenation (reduction)

Addition of hydrogen gas to alkene produces alkane with hydrogens in syn addition
Alkynes
- Deprotonation

To convert an alkyne to an alkyide anion, react it with a strong base such as an azanide
- Alkylation

Acetylide anions can replace halogens in alkanes to add an alkyl group to form an alkyne
- Alkyne From Alkene

Alkenes can form alkynes by reacting first with halogen gas to form a dihalogen alkane and then with some sodium amide and ammonia to form the alkyne
The second step removes 2 hydrogens and 2 halogens but can be tuned to remove only one
- Halogenation

Addition of one mole of halogen gives two halogen molecules added anti to each other in an alkene. Another molar addition of halogen gives an alkane saturated with the halogen. Carbocation rearrangment is not possible
- Hydrohalogenation

A molar addition of hydrogen halide produces an alkene with the halogen added to the most substitude carbon with possible carbocation rearrangmenets. Another molar addition of hydrogen halide produces an alkane that again follows Markovinikov’s rule

- Hydroboration

In the addition of borane with hydrogen peroxide and sodium hydroxide to an internal alkyne, a ketone is formed on either carbon. If the reaction is done on a terminal alkyne, the result is an alkane with a diol on the terminal carbon

In the addition of $\ce{(sia)2BH}$ with hydrogen peroxide and sodium hydroxide to a terminal alkyne, an aldehyde is formed at the very end of the alkane. Only one hydroboration happenes because (sia)2 borane is a weaker version of borane
- Hydration

In the addition of mercury sulfate with hydrosulfuric acid to an alkyne, a ketone is formed with the oxygen double bonded to the more substituted carbon without rearrangements

- Hydrogenation

In the addition of two moles of hydrogen gas to alkyne with a transition metal catalyst, an alkane is formed

In the addition of hydrogen gas to alkyne with Lindlar’s catalyst (transition metal poisoned with lead), two stereoisomers of cis-alkenes may be formed

In the addition of a carboxylic acid to a hydroboranated alkyne, a cis-alkene is produced along with a boron bonded to three carboxylate groups

In the addition of alkali metal to alkyne in liquid ammonia, an alkene is formed with the hydrogens added anti to each other

Haloalkanes
- Halogenation

In the addition of halogen to alkane, the halogen atom typically adds to the more substituted carbon in a radical chain reaction

- Allylic Halogenation

In the addition of NBS with carbon tetrachloride and light to alkene, 5

- Autooxidation

- Hydrohalogenation


| | $3^\degree:2^\degree:1^\degree$ |
| | |
| $\ce{Br2}$ | $1600:80:1$ |
| $\ce{Cl2}$ | $5:4:1$ |
Biology 🔬
Nervous System
Vm=zFRTln(cico)
Is one form of the Nernst Equation
A cell produces a voltage difference (inside - outside) with a separation of charge at a certain temperature T, charge carrier charge z, inside concentrations of ions ci, and outside concentrations of ions co
Vm=FRTln(PA[A]o+…+PN[N]oPA[A]i+…+PN[N]i) PA represents the permeability of the membrane to A.
Action Potential
CNS:
Brain/spinal
PNS:
Somatic
Control of skeletal muscles
Autonomic
Fight/flight (sympathetic)
Food/digestion (parasympathetic)
Frontal
Primary motor cortex, prefrontal cortex, Broca’s area (grammar wrong, meaning there)
Parietal lobe
Somatosensory cortex
Temporal
Auditory cortex, Wernicke’s area (word salad)
Occipital lobe
Primary visual cortex
Contralateral control, corpus callosem in between hemispheres
Thalaus
Sensory relay station
…
Quiz:
Amygdala involuntary movement
EEG on cerebral cortex
Cells
Cell Juncitons
Q=RP Q is flow rate, P is the pressure of the fluid, and R is the resistance to flow